update

* lora

* update

* update

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

.github/ISSUE_TEMPLATE/remote-vae-pilot-feedback.yml

@@ -0,0 +1,38 @@
+name: "\U0001F31F Remote VAE"
+description: Feedback for remote VAE pilot
+labels: [ "Remote VAE" ]
+
+body:
+  - type: textarea
+    id: positive
+    validations:
+      required: true
+    attributes:
+      label: Did you like the remote VAE solution?
+      description: |
+        If you liked it, we would appreciate it if you could elaborate what you liked.
+
+  - type: textarea
+    id: feedback
+    validations:
+      required: true
+    attributes:
+      label: What can be improved about the current solution?
+      description: |
+        Let us know the things you would like to see improved. Note that we will work optimizing the solution once the pilot is over and we have usage.
+
+  - type: textarea
+    id: others
+    validations:
+      required: true
+    attributes:
+      label: What other VAEs you would like to see if the pilot goes well?
+      description: |
+        Provide a list of the VAEs you would like to see in the future if the pilot goes well.
+
+  - type: textarea
+    id: additional-info
+    attributes:
+      label: Notify the members of the team
+      description: |
+        Tag the following folks when submitting this feedback: @hlky @sayakpaul

.github/workflows/nightly_tests.yml

@@ -265,7 +265,7 @@ jobs:
 
       - name: Run PyTorch CUDA tests
         env:
-          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
           # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
           CUBLAS_WORKSPACE_CONFIG: :16:8
         run: |
@@ -505,7 +505,7 @@ jobs:
 #        shell: arch -arch arm64 bash {0}
 #        env:
 #          HF_HOME: /System/Volumes/Data/mnt/cache
-#          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+#          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
 #        run: |
 #          ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
 #            --report-log=tests_torch_mps.log \
@@ -561,7 +561,7 @@ jobs:
 #        shell: arch -arch arm64 bash {0}
 #        env:
 #          HF_HOME: /System/Volumes/Data/mnt/cache
-#          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+#          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
 #        run: |
 #          ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
 #            --report-log=tests_torch_mps.log \

.github/workflows/pr_style_bot.yml

@@ -0,0 +1,127 @@
+name: PR Style Bot
+
+on:
+  issue_comment:
+    types: [created]
+
+permissions:
+  contents: write
+  pull-requests: write
+
+jobs:
+  run-style-bot:
+    if: >
+      contains(github.event.comment.body, '@bot /style') &&
+      github.event.issue.pull_request != null
+    runs-on: ubuntu-latest
+
+    steps:
+      - name: Extract PR details
+        id: pr_info
+        uses: actions/github-script@v6
+        with:
+          script: |
+            const prNumber = context.payload.issue.number;
+            const { data: pr } = await github.rest.pulls.get({
+              owner: context.repo.owner,
+              repo: context.repo.repo,
+              pull_number: prNumber
+            });
+            
+            // We capture both the branch ref and the "full_name" of the head repo
+            // so that we can check out the correct repository & branch (including forks).
+            core.setOutput("prNumber", prNumber);
+            core.setOutput("headRef", pr.head.ref);
+            core.setOutput("headRepoFullName", pr.head.repo.full_name);
+
+      - name: Check out PR branch
+        uses: actions/checkout@v3
+        env: 
+          HEADREPOFULLNAME: ${{ steps.pr_info.outputs.headRepoFullName }}
+          HEADREF: ${{ steps.pr_info.outputs.headRef }}
+        with:
+          # Instead of checking out the base repo, use the contributor's repo name
+          repository: ${{ env.HEADREPOFULLNAME }}
+          ref: ${{ env.HEADREF }}
+          # You may need fetch-depth: 0 for being able to push
+          fetch-depth: 0
+          token: ${{ secrets.GITHUB_TOKEN }}
+      
+      - name: Debug
+        env: 
+          HEADREPOFULLNAME: ${{ steps.pr_info.outputs.headRepoFullName }}
+          HEADREF: ${{ steps.pr_info.outputs.headRef }}
+          PRNUMBER: ${{ steps.pr_info.outputs.prNumber }}
+        run: |
+          echo "PR number: ${{ env.PRNUMBER }}"
+          echo "Head Ref: ${{ env.HEADREF }}"
+          echo "Head Repo Full Name: ${{ env.HEADREPOFULLNAME }}"
+
+      - name: Set up Python
+        uses: actions/setup-python@v4
+
+      - name: Install dependencies
+        run: |
+          pip install .[quality]
+
+      - name: Download Makefile from main branch
+        run: |
+          curl -o main_Makefile https://raw.githubusercontent.com/huggingface/diffusers/main/Makefile
+        
+      - name: Compare Makefiles
+        run: |
+          if ! diff -q main_Makefile Makefile; then
+            echo "Error: The Makefile has changed. Please ensure it matches the main branch."
+            exit 1
+          fi
+          echo "No changes in Makefile. Proceeding..."
+          rm -rf main_Makefile
+
+      - name: Run make style and make quality
+        run: |
+          make style && make quality
+
+      - name: Commit and push changes
+        id: commit_and_push
+        env: 
+          HEADREPOFULLNAME: ${{ steps.pr_info.outputs.headRepoFullName }}
+          HEADREF: ${{ steps.pr_info.outputs.headRef }}
+          PRNUMBER: ${{ steps.pr_info.outputs.prNumber }}
+          GITHUB_TOKEN: ${{ secrets.GITHUB_TOKEN }}
+        run: |
+          echo "HEADREPOFULLNAME: ${{ env.HEADREPOFULLNAME }}, HEADREF: ${{ env.HEADREF }}"
+          # Configure git with the Actions bot user
+          git config user.name "github-actions[bot]"
+          git config user.email "github-actions[bot]@users.noreply.github.com"
+
+          # Make sure your 'origin' remote is set to the contributor's fork
+          git remote set-url origin "https://x-access-token:${GITHUB_TOKEN}@github.com/${{ env.HEADREPOFULLNAME }}.git"
+
+          # If there are changes after running style/quality, commit them
+          if [ -n "$(git status --porcelain)" ]; then
+            git add .
+            git commit -m "Apply style fixes"
+            # Push to the original contributor's forked branch
+            git push origin HEAD:${{ env.HEADREF }}
+            echo "changes_pushed=true" >> $GITHUB_OUTPUT
+          else
+            echo "No changes to commit."
+            echo "changes_pushed=false" >> $GITHUB_OUTPUT
+          fi
+
+      - name: Comment on PR with workflow run link
+        if: steps.commit_and_push.outputs.changes_pushed == 'true'
+        uses: actions/github-script@v6
+        with:
+          script: |
+            const prNumber = parseInt(process.env.prNumber, 10);
+            const runUrl = `${process.env.GITHUB_SERVER_URL}/${process.env.GITHUB_REPOSITORY}/actions/runs/${process.env.GITHUB_RUN_ID}`
+
+            await github.rest.issues.createComment({
+              owner: context.repo.owner,
+              repo: context.repo.repo,
+              issue_number: prNumber,
+              body: `Style fixes have been applied. [View the workflow run here](${runUrl}).`
+            });
+        env:
+          prNumber: ${{ steps.pr_info.outputs.prNumber }}

.github/workflows/pr_tests.yml

@@ -2,8 +2,8 @@ name: Fast tests for PRs
 
 on:
   pull_request:
-    branches:
-      - main
+    branches: [main]
+    types: [synchronize]
     paths:
       - "src/diffusers/**.py"
       - "benchmarks/**.py"
@@ -64,6 +64,7 @@ jobs:
         run: |
           python utils/check_copies.py
           python utils/check_dummies.py
+          python utils/check_support_list.py
           make deps_table_check_updated
       - name: Check if failure
         if: ${{ failure() }}
@@ -120,7 +121,8 @@ jobs:
       run: |
         python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
         python -m uv pip install -e [quality,test]
-        python -m uv pip install accelerate
+        pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
+        pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
 
     - name: Environment
       run: |

.github/workflows/push_tests.yml

@@ -1,6 +1,13 @@
 name: Fast GPU Tests on main
 
 on:
+  pull_request:
+    branches: main
+    paths:
+      - "src/diffusers/models/modeling_utils.py"
+      - "src/diffusers/models/model_loading_utils.py"
+      - "src/diffusers/pipelines/pipeline_utils.py"
+      - "src/diffusers/pipeline_loading_utils.py"
   workflow_dispatch:
   push:
     branches:
@@ -83,7 +90,7 @@ jobs:
           python utils/print_env.py
       - name: PyTorch CUDA checkpoint tests on Ubuntu
         env:
-          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
           # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
           CUBLAS_WORKSPACE_CONFIG: :16:8
         run: |
@@ -137,7 +144,7 @@ jobs:
 
     - name: Run PyTorch CUDA tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
         # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
         CUBLAS_WORKSPACE_CONFIG: :16:8
       run: |
@@ -160,6 +167,7 @@ jobs:
         path: reports
 
   flax_tpu_tests:
+    if: ${{ github.event_name != 'pull_request' }}
     name: Flax TPU Tests
     runs-on:
       group: gcp-ct5lp-hightpu-8t
@@ -187,7 +195,7 @@ jobs:
 
     - name: Run Flax TPU tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 0 \
           -s -v -k "Flax" \
@@ -208,6 +216,7 @@ jobs:
         path: reports
 
   onnx_cuda_tests:
+    if: ${{ github.event_name != 'pull_request' }}
     name: ONNX CUDA Tests
     runs-on:
       group: aws-g4dn-2xlarge
@@ -235,7 +244,7 @@ jobs:
 
     - name: Run ONNXRuntime CUDA tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
           -s -v -k "Onnx" \
@@ -256,6 +265,7 @@ jobs:
         path: reports
 
   run_torch_compile_tests:
+    if: ${{ github.event_name != 'pull_request' }}
     name: PyTorch Compile CUDA tests
 
     runs-on:
@@ -283,7 +293,7 @@ jobs:
         python utils/print_env.py
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
         RUN_COMPILE: yes
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
@@ -299,6 +309,7 @@ jobs:
         path: reports
 
   run_xformers_tests:
+    if: ${{ github.event_name != 'pull_request' }}
     name: PyTorch xformers CUDA tests
 
     runs-on:
@@ -326,7 +337,7 @@ jobs:
         python utils/print_env.py
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
     - name: Failure short reports
@@ -349,7 +360,6 @@ jobs:
     container:
       image: diffusers/diffusers-pytorch-cuda
       options: --gpus 0 --shm-size "16gb" --ipc host
-
     steps:
     - name: Checkout diffusers
       uses: actions/checkout@v3
@@ -359,7 +369,6 @@ jobs:
     - name: NVIDIA-SMI
       run: |
         nvidia-smi
-
     - name: Install dependencies
       run: |
         python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
@@ -372,7 +381,7 @@ jobs:
 
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
         python -m uv pip install timm

.github/workflows/release_tests_fast.yml

@@ -81,7 +81,7 @@ jobs:
           python utils/print_env.py
       - name: Slow PyTorch CUDA checkpoint tests on Ubuntu
         env:
-          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
           # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
           CUBLAS_WORKSPACE_CONFIG: :16:8
         run: |
@@ -135,7 +135,7 @@ jobs:
 
     - name: Run PyTorch CUDA tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
         # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
         CUBLAS_WORKSPACE_CONFIG: :16:8
       run: |
@@ -186,7 +186,7 @@ jobs:
 
       - name: Run PyTorch CUDA tests
         env:
-          HF_TOKEN: ${{ secrets.HF_TOKEN }}
+          HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
           # https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
           CUBLAS_WORKSPACE_CONFIG: :16:8
         run: |
@@ -241,7 +241,7 @@ jobs:
 
     - name: Run slow Flax TPU tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 0 \
           -s -v -k "Flax" \
@@ -289,7 +289,7 @@ jobs:
 
     - name: Run slow ONNXRuntime CUDA tests
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
           -s -v -k "Onnx" \
@@ -337,7 +337,7 @@ jobs:
         python utils/print_env.py
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
         RUN_COMPILE: yes
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
@@ -380,7 +380,7 @@ jobs:
         python utils/print_env.py
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
     - name: Failure short reports
@@ -426,7 +426,7 @@ jobs:
 
     - name: Run example tests on GPU
       env:
-        HF_TOKEN: ${{ secrets.HF_TOKEN }}
+        HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
       run: |
         python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
         python -m uv pip install timm

.github/workflows/run_tests_from_a_pr.yml

@@ -7,8 +7,8 @@ on:
         default: 'diffusers/diffusers-pytorch-cuda'
         description: 'Name of the Docker image'
         required: true
-      branch:
-        description: 'PR Branch to test on'
+      pr_number:
+        description: 'PR number to test on'
         required: true
       test:
         description: 'Tests to run (e.g.: `tests/models`).'
@@ -43,8 +43,8 @@ jobs:
             exit 1
           fi
 
-          if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
-            echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
+          if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines|lora) ]]; then
+            echo "Error: The input string must contain either 'models', 'pipelines', or 'lora' after 'tests/'."
             exit 1
           fi
 
@@ -53,13 +53,13 @@ jobs:
             exit 1
           fi
           echo "$PY_TEST"
+        
+        shell: bash -e {0}
 
       - name: Checkout PR branch
         uses: actions/checkout@v4
         with:
-          ref: ${{ github.event.inputs.branch }}
-          repository: ${{ github.event.pull_request.head.repo.full_name }}
-
+          ref: refs/pull/${{ inputs.pr_number }}/head
 
       - name: Install pytest
         run: |

.github/workflows/trufflehog.yml

@@ -13,3 +13,6 @@ jobs:
         fetch-depth: 0
     - name: Secret Scanning
       uses: trufflesecurity/trufflehog@main
+      with:
+        extra_args: --results=verified,unknown
+

docs/source/en/_toctree.yml

@@ -79,6 +79,8 @@
 - sections:
   - local: using-diffusers/cogvideox
     title: CogVideoX
+  - local: using-diffusers/consisid
+    title: ConsisID
   - local: using-diffusers/sdxl
     title: Stable Diffusion XL
   - local: using-diffusers/sdxl_turbo
@@ -87,6 +89,8 @@
     title: Kandinsky
   - local: using-diffusers/ip_adapter
     title: IP-Adapter
+  - local: using-diffusers/omnigen
+    title: OmniGen
   - local: using-diffusers/pag
     title: PAG
   - local: using-diffusers/controlnet
@@ -179,6 +183,8 @@
     title: TGATE
   - local: optimization/xdit
     title: xDiT
+  - local: optimization/para_attn
+    title: ParaAttention
   - sections:
     - local: using-diffusers/stable_diffusion_jax_how_to
       title: JAX/Flax
@@ -268,8 +274,12 @@
         title: AuraFlowTransformer2DModel
       - local: api/models/cogvideox_transformer3d
         title: CogVideoXTransformer3DModel
+      - local: api/models/consisid_transformer3d
+        title: ConsisIDTransformer3DModel
       - local: api/models/cogview3plus_transformer2d
         title: CogView3PlusTransformer2DModel
+      - local: api/models/cogview4_transformer2d
+        title: CogView4Transformer2DModel
       - local: api/models/dit_transformer2d
         title: DiTTransformer2DModel
       - local: api/models/flux_transformer
@@ -282,10 +292,14 @@
         title: LatteTransformer3DModel
       - local: api/models/lumina_nextdit2d
         title: LuminaNextDiT2DModel
+      - local: api/models/lumina2_transformer2d
+        title: Lumina2Transformer2DModel
       - local: api/models/ltx_video_transformer3d
         title: LTXVideoTransformer3DModel
       - local: api/models/mochi_transformer3d
         title: MochiTransformer3DModel
+      - local: api/models/omnigen_transformer
+        title: OmniGenTransformer2DModel
       - local: api/models/pixart_transformer2d
         title: PixArtTransformer2DModel
       - local: api/models/prior_transformer
@@ -370,6 +384,10 @@
       title: CogVideoX
     - local: api/pipelines/cogview3
       title: CogView3
+    - local: api/pipelines/cogview4
+      title: CogView4
+    - local: api/pipelines/consisid
+      title: ConsisID
     - local: api/pipelines/consistency_models
       title: Consistency Models
     - local: api/pipelines/controlnet
@@ -430,6 +448,8 @@
       title: LEDITS++
     - local: api/pipelines/ltx_video
       title: LTXVideo
+    - local: api/pipelines/lumina2
+      title: Lumina 2.0
     - local: api/pipelines/lumina
       title: Lumina-T2X
     - local: api/pipelines/marigold
@@ -440,6 +460,8 @@
       title: MultiDiffusion
     - local: api/pipelines/musicldm
       title: MusicLDM
+    - local: api/pipelines/omnigen
+      title: OmniGen
     - local: api/pipelines/pag
       title: PAG
     - local: api/pipelines/paint_by_example
@@ -590,6 +612,8 @@
       title: Attention Processor
     - local: api/activations
       title: Custom activation functions
+    - local: api/cache
+      title: Caching methods
     - local: api/normalization
       title: Custom normalization layers
     - local: api/utilities

docs/source/en/api/activations.md

@@ -25,3 +25,16 @@ Customized activation functions for supporting various models in 🤗 Diffusers.
 ## ApproximateGELU
 
 [[autodoc]] models.activations.ApproximateGELU
+
+
+## SwiGLU
+
+[[autodoc]] models.activations.SwiGLU
+
+## FP32SiLU
+
+[[autodoc]] models.activations.FP32SiLU
+
+## LinearActivation
+
+[[autodoc]] models.activations.LinearActivation

docs/source/en/api/attnprocessor.md

@@ -147,3 +147,20 @@ An attention processor is a class for applying different types of attention mech
 ## XLAFlashAttnProcessor2_0
 
 [[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
+
+## XFormersJointAttnProcessor
+
+[[autodoc]] models.attention_processor.XFormersJointAttnProcessor
+
+## IPAdapterXFormersAttnProcessor
+
+[[autodoc]] models.attention_processor.IPAdapterXFormersAttnProcessor
+
+## FluxIPAdapterJointAttnProcessor2_0
+
+[[autodoc]] models.attention_processor.FluxIPAdapterJointAttnProcessor2_0
+
+
+## XLAFluxFlashAttnProcessor2_0
+
+[[autodoc]] models.attention_processor.XLAFluxFlashAttnProcessor2_0

docs/source/en/api/cache.md

@@ -0,0 +1,49 @@
+<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License. -->
+
+# Caching methods
+
+## Pyramid Attention Broadcast
+
+[Pyramid Attention Broadcast](https://huggingface.co/papers/2408.12588) from Xuanlei Zhao, Xiaolong Jin, Kai Wang, Yang You.
+
+Pyramid Attention Broadcast (PAB) is a method that speeds up inference in diffusion models by systematically skipping attention computations between successive inference steps and reusing cached attention states. The attention states are not very different between successive inference steps. The most prominent difference is in the spatial attention blocks, not as much in the temporal attention blocks, and finally the least in the cross attention blocks. Therefore, many cross attention computation blocks can be skipped, followed by the temporal and spatial attention blocks. By combining other techniques like sequence parallelism and classifier-free guidance parallelism, PAB achieves near real-time video generation.
+
+Enable PAB with [`~PyramidAttentionBroadcastConfig`] on any pipeline. For some benchmarks, refer to [this](https://github.com/huggingface/diffusers/pull/9562) pull request.
+
+```python
+import torch
+from diffusers import CogVideoXPipeline, PyramidAttentionBroadcastConfig
+
+pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
+pipe.to("cuda")
+
+# Increasing the value of `spatial_attention_timestep_skip_range[0]` or decreasing the value of
+# `spatial_attention_timestep_skip_range[1]` will decrease the interval in which pyramid attention
+# broadcast is active, leader to slower inference speeds. However, large intervals can lead to
+# poorer quality of generated videos.
+config = PyramidAttentionBroadcastConfig(
+    spatial_attention_block_skip_range=2,
+    spatial_attention_timestep_skip_range=(100, 800),
+    current_timestep_callback=lambda: pipe.current_timestep,
+)
+pipe.transformer.enable_cache(config)
+```
+
+### CacheMixin
+
+[[autodoc]] CacheMixin
+
+### PyramidAttentionBroadcastConfig
+
+[[autodoc]] PyramidAttentionBroadcastConfig
+
+[[autodoc]] apply_pyramid_attention_broadcast

docs/source/en/api/loaders/lora.md

@@ -20,6 +20,10 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
 - [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
 - [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
 - [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
+- [`LTXVideoLoraLoaderMixin`] provides similar functions for [LTX-Video](https://huggingface.co/docs/diffusers/main/en/api/pipelines/ltx_video).
+- [`SanaLoraLoaderMixin`] provides similar functions for [Sana](https://huggingface.co/docs/diffusers/main/en/api/pipelines/sana).
+- [`HunyuanVideoLoraLoaderMixin`] provides similar functions for [HunyuanVideo](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hunyuan_video).
+- [`Lumina2LoraLoaderMixin`] provides similar functions for [Lumina2](https://huggingface.co/docs/diffusers/main/en/api/pipelines/lumina2).
 - [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
 - [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
 
@@ -53,6 +57,22 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
 
 [[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
 
+## LTXVideoLoraLoaderMixin
+
+[[autodoc]] loaders.lora_pipeline.LTXVideoLoraLoaderMixin
+
+## SanaLoraLoaderMixin
+
+[[autodoc]] loaders.lora_pipeline.SanaLoraLoaderMixin
+
+## HunyuanVideoLoraLoaderMixin
+
+[[autodoc]] loaders.lora_pipeline.HunyuanVideoLoraLoaderMixin
+
+## Lumina2LoraLoaderMixin
+
+[[autodoc]] loaders.lora_pipeline.Lumina2LoraLoaderMixin
+
 ## AmusedLoraLoaderMixin
 
 [[autodoc]] loaders.lora_pipeline.AmusedLoraLoaderMixin

docs/source/en/api/models/cogview4_transformer2d.md

@@ -0,0 +1,30 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License. -->
+
+# CogView4Transformer2DModel
+
+A Diffusion Transformer model for 2D data from [CogView4]()
+
+The model can be loaded with the following code snippet.
+
+```python
+from diffusers import CogView4Transformer2DModel
+
+transformer = CogView4Transformer2DModel.from_pretrained("THUDM/CogView4-6B", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
+```
+
+## CogView4Transformer2DModel
+
+[[autodoc]] CogView4Transformer2DModel
+
+## Transformer2DModelOutput
+
+[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

docs/source/en/api/models/consisid_transformer3d.md

@@ -0,0 +1,30 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License. -->
+
+# ConsisIDTransformer3DModel
+
+A Diffusion Transformer model for 3D data from [ConsisID](https://github.com/PKU-YuanGroup/ConsisID) was introduced in [Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/pdf/2411.17440) by Peking University & University of Rochester & etc.
+
+The model can be loaded with the following code snippet.
+
+```python
+from diffusers import ConsisIDTransformer3DModel
+
+transformer = ConsisIDTransformer3DModel.from_pretrained("BestWishYsh/ConsisID-preview", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
+```
+
+## ConsisIDTransformer3DModel
+
+[[autodoc]] ConsisIDTransformer3DModel
+
+## Transformer2DModelOutput
+
+[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

docs/source/en/api/models/lumina2_transformer2d.md

@@ -0,0 +1,30 @@
+<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License. -->
+
+# Lumina2Transformer2DModel
+
+A Diffusion Transformer model for 3D video-like data was introduced in [Lumina Image 2.0](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) by Alpha-VLLM.
+
+The model can be loaded with the following code snippet.
+
+```python
+from diffusers import Lumina2Transformer2DModel
+
+transformer = Lumina2Transformer2DModel.from_pretrained("Alpha-VLLM/Lumina-Image-2.0", subfolder="transformer", torch_dtype=torch.bfloat16)
+```
+
+## Lumina2Transformer2DModel
+
+[[autodoc]] Lumina2Transformer2DModel
+
+## Transformer2DModelOutput
+
+[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

docs/source/en/api/models/omnigen_transformer.md

@@ -0,0 +1,30 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License.
+-->
+
+# OmniGenTransformer2DModel
+
+A Transformer model that accepts multimodal instructions to generate images for [OmniGen](https://github.com/VectorSpaceLab/OmniGen/).
+
+The abstract from the paper is:
+
+*The emergence of Large Language Models (LLMs) has unified language  generation tasks and revolutionized human-machine interaction.  However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the model’s reasoning capabilities and potential applications of the chain-of-thought mechanism.  This work represents the first attempt at a general-purpose image generation model,  and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
+
+```python
+import torch
+from diffusers import OmniGenTransformer2DModel
+
+transformer = OmniGenTransformer2DModel.from_pretrained("Shitao/OmniGen-v1-diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
+```
+
+## OmniGenTransformer2DModel
+
+[[autodoc]] OmniGenTransformer2DModel

docs/source/en/api/normalization.md

@@ -29,3 +29,43 @@ Customized normalization layers for supporting various models in 🤗 Diffusers.
 ## AdaGroupNorm
 
 [[autodoc]] models.normalization.AdaGroupNorm
+
+## AdaLayerNormContinuous
+
+[[autodoc]] models.normalization.AdaLayerNormContinuous
+
+## RMSNorm
+
+[[autodoc]] models.normalization.RMSNorm
+
+## GlobalResponseNorm
+
+[[autodoc]] models.normalization.GlobalResponseNorm
+
+
+## LuminaLayerNormContinuous
+[[autodoc]] models.normalization.LuminaLayerNormContinuous
+
+## SD35AdaLayerNormZeroX
+[[autodoc]] models.normalization.SD35AdaLayerNormZeroX
+
+## AdaLayerNormZeroSingle
+[[autodoc]] models.normalization.AdaLayerNormZeroSingle
+
+## LuminaRMSNormZero
+[[autodoc]] models.normalization.LuminaRMSNormZero
+
+## LpNorm
+[[autodoc]] models.normalization.LpNorm
+
+## CogView3PlusAdaLayerNormZeroTextImage
+[[autodoc]] models.normalization.CogView3PlusAdaLayerNormZeroTextImage
+
+## CogVideoXLayerNormZero
+[[autodoc]] models.normalization.CogVideoXLayerNormZero
+
+## MochiRMSNormZero
+[[autodoc]] models.transformers.transformer_mochi.MochiRMSNormZero
+
+## MochiRMSNorm
+[[autodoc]] models.normalization.MochiRMSNorm

docs/source/en/api/pipelines/animatediff.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Text-to-Video Generation with AnimateDiff
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ## Overview
 
 [AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.

docs/source/en/api/pipelines/cogvideox.md

@@ -15,6 +15,10 @@
 
 # CogVideoX
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/cogview4.md

@@ -0,0 +1,34 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+#
+# Licensed under the Apache License, Version 2.0 (the "License");
+# you may not use this file except in compliance with the License.
+# You may obtain a copy of the License at
+#
+#     http://www.apache.org/licenses/LICENSE-2.0
+#
+# Unless required by applicable law or agreed to in writing, software
+# distributed under the License is distributed on an "AS IS" BASIS,
+# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
+# See the License for the specific language governing permissions and
+# limitations under the License.
+-->
+
+# CogView4
+
+<Tip>
+
+Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
+
+</Tip>
+
+This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
+
+## CogView4Pipeline
+
+[[autodoc]] CogView4Pipeline
+  - all
+  - __call__
+
+## CogView4PipelineOutput
+
+[[autodoc]] pipelines.cogview4.pipeline_output.CogView4PipelineOutput

docs/source/en/api/pipelines/consisid.md

@@ -0,0 +1,64 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+#
+# Licensed under the Apache License, Version 2.0 (the "License");
+# you may not use this file except in compliance with the License.
+# You may obtain a copy of the License at
+#
+#     http://www.apache.org/licenses/LICENSE-2.0
+#
+# Unless required by applicable law or agreed to in writing, software
+# distributed under the License is distributed on an "AS IS" BASIS,
+# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
+# See the License for the specific language governing permissions and
+# limitations under the License.
+-->
+
+# ConsisID
+
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
+[Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/abs/2411.17440) from Peking University & University of Rochester & etc, by Shenghai Yuan, Jinfa Huang, Xianyi He, Yunyang Ge, Yujun Shi, Liuhan Chen, Jiebo Luo, Li Yuan.
+
+The abstract from the paper is:
+
+*Identity-preserving text-to-video (IPT2V) generation aims to create high-fidelity videos with consistent human identity. It is an important task in video generation but remains an open problem for generative models. This paper pushes the technical frontier of IPT2V in two directions that have not been resolved in the literature: (1) A tuning-free pipeline without tedious case-by-case finetuning, and (2) A frequency-aware heuristic identity-preserving Diffusion Transformer (DiT)-based control scheme. To achieve these goals, we propose **ConsisID**, a tuning-free DiT-based controllable IPT2V model to keep human-**id**entity **consis**tent in the generated video. Inspired by prior findings in frequency analysis of vision/diffusion transformers, it employs identity-control signals in the frequency domain, where facial features can be decomposed into low-frequency global features (e.g., profile, proportions) and high-frequency intrinsic features (e.g., identity markers that remain unaffected by pose changes). First, from a low-frequency perspective, we introduce a global facial extractor, which encodes the reference image and facial key points into a latent space, generating features enriched with low-frequency information. These features are then integrated into the shallow layers of the network to alleviate training challenges associated with DiT. Second, from a high-frequency perspective, we design a local facial extractor to capture high-frequency details and inject them into the transformer blocks, enhancing the model's ability to preserve fine-grained features. To leverage the frequency information for identity preservation, we propose a hierarchical training strategy, transforming a vanilla pre-trained video generation model into an IPT2V model. Extensive experiments demonstrate that our frequency-aware heuristic scheme provides an optimal control solution for DiT-based models. Thanks to this scheme, our **ConsisID** achieves excellent results in generating high-quality, identity-preserving videos, making strides towards more effective IPT2V. The model weight of ConsID is publicly available at https://github.com/PKU-YuanGroup/ConsisID.*
+
+<Tip>
+
+Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
+
+</Tip>
+
+This pipeline was contributed by [SHYuanBest](https://github.com/SHYuanBest). The original codebase can be found [here](https://github.com/PKU-YuanGroup/ConsisID). The original weights can be found under [hf.co/BestWishYsh](https://huggingface.co/BestWishYsh).
+
+There are two official ConsisID checkpoints for identity-preserving text-to-video.
+
+| checkpoints | recommended inference dtype |
+|:---:|:---:|
+| [`BestWishYsh/ConsisID-preview`](https://huggingface.co/BestWishYsh/ConsisID-preview) | torch.bfloat16 |
+| [`BestWishYsh/ConsisID-1.5`](https://huggingface.co/BestWishYsh/ConsisID-preview) | torch.bfloat16 |
+
+### Memory optimization
+
+ConsisID requires about 44 GB of GPU memory to decode 49 frames (6 seconds of video at 8 FPS) with output resolution 720x480 (W x H), which makes it not possible to run on consumer GPUs or free-tier T4 Colab. The following memory optimizations could be used to reduce the memory footprint. For replication, you can refer to [this](https://gist.github.com/SHYuanBest/bc4207c36f454f9e969adbb50eaf8258) script.
+
+| Feature (overlay the previous) | Max Memory Allocated | Max Memory Reserved |
+| :----------------------------- | :------------------- | :------------------ |
+| -                              | 37 GB                | 44 GB               |
+| enable_model_cpu_offload       | 22 GB                | 25 GB               |
+| enable_sequential_cpu_offload  | 16 GB                | 22 GB               |
+| vae.enable_slicing             | 16 GB                | 22 GB               |
+| vae.enable_tiling              | 5 GB                 | 7 GB                |
+
+## ConsisIDPipeline
+
+[[autodoc]] ConsisIDPipeline
+
+  - all
+  - __call__
+
+## ConsisIDPipelineOutput
+
+[[autodoc]] pipelines.consisid.pipeline_output.ConsisIDPipelineOutput

docs/source/en/api/pipelines/control_flux_inpaint.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # FluxControlInpaint
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 FluxControlInpaintPipeline is an implementation of Inpainting for Flux.1 Depth/Canny models. It is a pipeline that allows you to inpaint images using the Flux.1 Depth/Canny models. The pipeline takes an image and a mask as input and returns the inpainted image.
 
 FLUX.1 Depth and Canny [dev] is a 12 billion parameter rectified flow transformer capable of generating an image based on a text description while following the structure of a given input image. **This is not a ControlNet model**.

docs/source/en/api/pipelines/controlnet.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNet
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
 
 With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

docs/source/en/api/pipelines/controlnet_flux.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNet with Flux.1
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
 
 ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

docs/source/en/api/pipelines/controlnet_sd3.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNet with Stable Diffusion 3
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 StableDiffusion3ControlNetPipeline is an implementation of ControlNet for Stable Diffusion 3.
 
 ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

docs/source/en/api/pipelines/controlnet_sdxl.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNet with Stable Diffusion XL
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
 
 With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

docs/source/en/api/pipelines/controlnet_union.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNetUnion
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
 
 The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.

docs/source/en/api/pipelines/controlnetxs.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # ControlNet-XS
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
 
 Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

docs/source/en/api/pipelines/deepfloyd_if.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # DeepFloyd IF
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ## Overview
 
 DeepFloyd IF is a novel state-of-the-art open-source text-to-image model with a high degree of photorealism and language understanding.

docs/source/en/api/pipelines/flux.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Flux
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
 
 Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
@@ -309,6 +313,53 @@ image.save("output.png")
 
 When unloading the Control LoRA weights, call `pipe.unload_lora_weights(reset_to_overwritten_params=True)` to reset the `pipe.transformer` completely back to its original form. The resultant pipeline can then be used with methods like [`DiffusionPipeline.from_pipe`]. More details about this argument are available in [this PR](https://github.com/huggingface/diffusers/pull/10397).
 
+## IP-Adapter
+
+<Tip>
+
+Check out [IP-Adapter](../../../using-diffusers/ip_adapter) to learn more about how IP-Adapters work.
+
+</Tip>
+
+An IP-Adapter lets you prompt Flux with images, in addition to the text prompt. This is especially useful when describing complex concepts that are difficult to articulate through text alone and you have reference images.
+
+```python
+import torch
+from diffusers import FluxPipeline
+from diffusers.utils import load_image
+
+pipe = FluxPipeline.from_pretrained(
+    "black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16
+).to("cuda")
+
+image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_input.jpg").resize((1024, 1024))
+
+pipe.load_ip_adapter(
+    "XLabs-AI/flux-ip-adapter",
+    weight_name="ip_adapter.safetensors",
+    image_encoder_pretrained_model_name_or_path="openai/clip-vit-large-patch14"
+)
+pipe.set_ip_adapter_scale(1.0)
+
+image = pipe(
+    width=1024,
+    height=1024,
+    prompt="wearing sunglasses",
+    negative_prompt="",
+    true_cfg=4.0,
+    generator=torch.Generator().manual_seed(4444),
+    ip_adapter_image=image,
+).images[0]
+
+image.save('flux_ip_adapter_output.jpg')
+```
+
+<div class="justify-center">
+    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_output.jpg"/>
+    <figcaption class="mt-2 text-sm text-center text-gray-500">IP-Adapter examples with prompt "wearing sunglasses"</figcaption>
+</div>
+
+
 ## Running FP16 inference
 
 Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.

docs/source/en/api/pipelines/hunyuan_video.md

@@ -14,6 +14,10 @@
 
 # HunyuanVideo
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
 
 *Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/tencent/HunyuanVideo).*
@@ -32,6 +36,21 @@ Recommendations for inference:
 - For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
 - For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
 
+## Available models
+
+The following models are available for the [`HunyuanVideoPipeline`](text-to-video) pipeline:
+
+| Model name | Description |
+|:---|:---|
+| [`hunyuanvideo-community/HunyuanVideo`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo) | Official HunyuanVideo (guidance-distilled). Performs best at multiple resolutions and frames. Performs best with `guidance_scale=6.0`, `true_cfg_scale=1.0` and without a negative prompt. |
+| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
+
+The following models are available for the image-to-video pipeline:
+
+| Model name | Description |
+|:---|:---|
+| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution. Performs best at `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
+
 ## Quantization
 
 Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.

docs/source/en/api/pipelines/kandinsky3.md

@@ -9,6 +9,10 @@ specific language governing permissions and limitations under the License.
 
 # Kandinsky 3
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
 
 The description from it's GitHub page:

docs/source/en/api/pipelines/kolors.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Kolors: Effective Training of Diffusion Model for Photorealistic Text-to-Image Synthesis
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/kolors_header_collage.png)
 
 Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).

docs/source/en/api/pipelines/latent_consistency_models.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Latent Consistency Models
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://huggingface.co/papers/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
 
 The abstract of the paper is as follows:

docs/source/en/api/pipelines/ledits_pp.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # LEDITS++
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/ltx_video.md

@@ -14,6 +14,10 @@
 
 # LTX Video
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
 
 <Tip>

docs/source/en/api/pipelines/lumina2.md

@@ -0,0 +1,87 @@
+<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
+#
+# Licensed under the Apache License, Version 2.0 (the "License");
+# you may not use this file except in compliance with the License.
+# You may obtain a copy of the License at
+#
+#     http://www.apache.org/licenses/LICENSE-2.0
+#
+# Unless required by applicable law or agreed to in writing, software
+# distributed under the License is distributed on an "AS IS" BASIS,
+# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
+# See the License for the specific language governing permissions and
+# limitations under the License. -->
+
+# Lumina2
+
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
+[Lumina Image 2.0: A Unified and Efficient Image Generative Model](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) is a 2 billion parameter flow-based diffusion transformer capable of generating diverse images from text descriptions.
+
+The abstract from the paper is:
+
+*We introduce Lumina-Image 2.0, an advanced text-to-image model that surpasses previous state-of-the-art methods across multiple benchmarks, while also shedding light on its potential to evolve into a generalist vision intelligence model. Lumina-Image 2.0 exhibits three key properties: (1) Unification – it adopts a unified architecture that treats text and image tokens as a joint sequence, enabling natural cross-modal interactions and facilitating task expansion. Besides, since high-quality captioners can provide semantically better-aligned text-image training pairs, we introduce a unified captioning system, UniCaptioner, which generates comprehensive and precise captions for the model. This not only accelerates model convergence but also enhances prompt adherence, variable-length prompt handling, and task generalization via prompt templates. (2) Efficiency – to improve the efficiency of the unified architecture, we develop a set of optimization techniques that improve semantic learning and fine-grained texture generation during training while incorporating inference-time acceleration strategies without compromising image quality. (3) Transparency – we open-source all training details, code, and models to ensure full reproducibility, aiming to bridge the gap between well-resourced closed-source research teams and independent developers.*
+
+<Tip>
+
+Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
+
+</Tip>
+
+## Using Single File loading with Lumina Image 2.0
+
+Single file loading for Lumina Image 2.0 is available for the `Lumina2Transformer2DModel`
+
+```python
+import torch
+from diffusers import Lumina2Transformer2DModel, Lumina2Text2ImgPipeline
+
+ckpt_path = "https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0/blob/main/consolidated.00-of-01.pth"
+transformer = Lumina2Transformer2DModel.from_single_file(
+    ckpt_path, torch_dtype=torch.bfloat16
+)
+
+pipe = Lumina2Text2ImgPipeline.from_pretrained(
+    "Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
+)
+pipe.enable_model_cpu_offload()
+image = pipe(
+    "a cat holding a sign that says hello",
+    generator=torch.Generator("cpu").manual_seed(0),
+).images[0]
+image.save("lumina-single-file.png")
+
+```
+
+## Using GGUF Quantized Checkpoints with Lumina Image 2.0
+
+GGUF Quantized checkpoints for the `Lumina2Transformer2DModel` can be loaded via `from_single_file` with the `GGUFQuantizationConfig` 
+
+```python
+from diffusers import Lumina2Transformer2DModel, Lumina2Text2ImgPipeline, GGUFQuantizationConfig 
+
+ckpt_path = "https://huggingface.co/calcuis/lumina-gguf/blob/main/lumina2-q4_0.gguf"
+transformer = Lumina2Transformer2DModel.from_single_file(
+    ckpt_path,
+    quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
+    torch_dtype=torch.bfloat16,
+)
+
+pipe = Lumina2Text2ImgPipeline.from_pretrained(
+    "Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
+)
+pipe.enable_model_cpu_offload()
+image = pipe(
+    "a cat holding a sign that says hello",
+    generator=torch.Generator("cpu").manual_seed(0),
+).images[0]
+image.save("lumina-gguf.png")
+```
+
+## Lumina2Text2ImgPipeline
+
+[[autodoc]] Lumina2Text2ImgPipeline
+  - all
+  - __call__

docs/source/en/api/pipelines/mochi.md

@@ -15,6 +15,10 @@
 
 # Mochi 1 Preview
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 > [!TIP]
 > Only a research preview of the model weights is available at the moment.
 
@@ -115,7 +119,7 @@ export_to_video(frames, "mochi.mp4", fps=30)
 
 ## Reproducing the results from the Genmo Mochi repo
 
-The [Genmo Mochi implementation](https://github.com/genmoai/mochi/tree/main) uses different precision values for each stage in the inference process. The text encoder and VAE use `torch.float32`, while the DiT uses `torch.bfloat16` with the [attention kernel](https://pytorch.org/docs/stable/generated/torch.nn.attention.sdpa_kernel.html#torch.nn.attention.sdpa_kernel) set to `EFFICIENT_ATTENTION`. Diffusers pipelines currently do not support setting different `dtypes` for different stages of the pipeline. In order to run inference in the same way as the the original implementation, please refer to the following example.
+The [Genmo Mochi implementation](https://github.com/genmoai/mochi/tree/main) uses different precision values for each stage in the inference process. The text encoder and VAE use `torch.float32`, while the DiT uses `torch.bfloat16` with the [attention kernel](https://pytorch.org/docs/stable/generated/torch.nn.attention.sdpa_kernel.html#torch.nn.attention.sdpa_kernel) set to `EFFICIENT_ATTENTION`. Diffusers pipelines currently do not support setting different `dtypes` for different stages of the pipeline. In order to run inference in the same way as the original implementation, please refer to the following example.
 
 <Tip>
 The original Mochi implementation zeros out empty prompts. However, enabling this option and placing the entire pipeline under autocast can lead to numerical overflows with the T5 text encoder.

docs/source/en/api/pipelines/omnigen.md

@@ -0,0 +1,80 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+#
+# Licensed under the Apache License, Version 2.0 (the "License");
+# you may not use this file except in compliance with the License.
+# You may obtain a copy of the License at
+#
+#     http://www.apache.org/licenses/LICENSE-2.0
+#
+# Unless required by applicable law or agreed to in writing, software
+# distributed under the License is distributed on an "AS IS" BASIS,
+# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
+# See the License for the specific language governing permissions and
+# limitations under the License.
+-->
+
+# OmniGen
+
+[OmniGen: Unified Image Generation](https://arxiv.org/pdf/2409.11340) from BAAI, by Shitao Xiao, Yueze Wang, Junjie Zhou, Huaying Yuan, Xingrun Xing, Ruiran Yan, Chaofan Li, Shuting Wang, Tiejun Huang, Zheng Liu.
+
+The abstract from the paper is:
+
+*The emergence of Large Language Models (LLMs) has unified language  generation tasks and revolutionized human-machine interaction.  However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the model’s reasoning capabilities and potential applications of the chain-of-thought mechanism.  This work represents the first attempt at a general-purpose image generation model,  and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
+
+<Tip>
+
+Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
+
+</Tip>
+
+This pipeline was contributed by [staoxiao](https://github.com/staoxiao). The original codebase can be found [here](https://github.com/VectorSpaceLab/OmniGen). The original weights can be found under [hf.co/shitao](https://huggingface.co/Shitao/OmniGen-v1).
+
+## Inference
+
+First, load the pipeline:
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+
+pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
+pipe.to("cuda")
+```
+
+For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image. 
+You can try setting the `height` and `width` parameters to generate images with different size.
+
+```python
+prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
+image = pipe(
+    prompt=prompt,
+    height=1024,
+    width=1024,
+    guidance_scale=3,
+    generator=torch.Generator(device="cpu").manual_seed(111),
+).images[0]
+image.save("output.png")
+```
+
+OmniGen supports multimodal inputs. 
+When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image. 
+It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
+
+```python
+prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(222)).images[0]
+image.save("output.png")
+```
+
+## OmniGenPipeline
+
+[[autodoc]] OmniGenPipeline
+  - all
+  - __call__

docs/source/en/api/pipelines/overview.md

@@ -54,7 +54,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
 | [DiT](dit) | text2image |
 | [Flux](flux) | text2image |
 | [Hunyuan-DiT](hunyuandit) | text2image |
-| [I2VGen-XL](i2vgenxl) | text2video |
+| [I2VGen-XL](i2vgenxl) | image2video |
 | [InstructPix2Pix](pix2pix) | image editing |
 | [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
 | [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |

docs/source/en/api/pipelines/pag.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Perturbed-Attention Guidance
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [Perturbed-Attention Guidance (PAG)](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) is a new diffusion sampling guidance that improves sample quality across both unconditional and conditional settings, achieving this without requiring further training or the integration of external modules.
 
 PAG was introduced in [Self-Rectifying Diffusion Sampling with Perturbed-Attention Guidance](https://huggingface.co/papers/2403.17377) by Donghoon Ahn, Hyoungwon Cho, Jaewon Min, Wooseok Jang, Jungwoo Kim, SeonHwa Kim, Hyun Hee Park, Kyong Hwan Jin and Seungryong Kim.

docs/source/en/api/pipelines/panorama.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # MultiDiffusion
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/pia.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Image-to-Video Generation with PIA (Personalized Image Animator)
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 ## Overview
 
 [PIA: Your Personalized Image Animator via Plug-and-Play Modules in Text-to-Image Models](https://arxiv.org/abs/2312.13964) by Yiming Zhang, Zhening Xing, Yanhong Zeng, Youqing Fang, Kai Chen

docs/source/en/api/pipelines/pix2pix.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # InstructPix2Pix
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/papers/2211.09800) is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/sana.md

@@ -14,6 +14,10 @@
 
 # SanaPipeline
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/stable_diffusion/depth2img.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Depth-to-image
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The Stable Diffusion model can also infer depth based on an image using [MiDaS](https://github.com/isl-org/MiDaS). This allows you to pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the image structure.
 
 <Tip>

docs/source/en/api/pipelines/stable_diffusion/img2img.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Image-to-image
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
 
 The [`StableDiffusionImg2ImgPipeline`] uses the diffusion-denoising mechanism proposed in [SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations](https://huggingface.co/papers/2108.01073) by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.

docs/source/en/api/pipelines/stable_diffusion/inpaint.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Inpainting
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The Stable Diffusion model can also be applied to inpainting which lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
 
 ## Tips

docs/source/en/api/pipelines/stable_diffusion/ldm3d_diffusion.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Text-to-(RGB, depth)
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
 
 Two checkpoints are available for use:

docs/source/en/api/pipelines/stable_diffusion/overview.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Stable Diffusion pipelines
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. This specific type of diffusion model was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
 
 Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.

docs/source/en/api/pipelines/stable_diffusion/stable_diffusion_3.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Stable Diffusion 3
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/stable_diffusion/stable_diffusion_xl.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Stable Diffusion XL
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/stable_diffusion/text2img.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Text-to-image
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The Stable Diffusion model was created by researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [Runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/stable_diffusion/upscale.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Super-resolution
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The Stable Diffusion upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/). It is used to enhance the resolution of input images by a factor of 4.
 
 <Tip>

docs/source/en/api/pipelines/stable_unclip.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Stable unCLIP
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 Stable unCLIP checkpoints are finetuned from [Stable Diffusion 2.1](./stable_diffusion/stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
 Stable unCLIP still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
 for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.

docs/source/en/api/pipelines/text_to_video.md

@@ -18,6 +18,10 @@ specific language governing permissions and limitations under the License.
 
 # Text-to-video
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [ModelScope Text-to-Video Technical Report](https://arxiv.org/abs/2308.06571) is by Jiuniu Wang, Hangjie Yuan, Dayou Chen, Yingya Zhang, Xiang Wang, Shiwei Zhang.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/text_to_video_zero.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Text2Video-Zero
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 [Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://huggingface.co/papers/2303.13439) is by Levon Khachatryan, Andranik Movsisyan, Vahram Tadevosyan, Roberto Henschel, [Zhangyang Wang](https://www.ece.utexas.edu/people/faculty/atlas-wang), Shant Navasardyan, [Humphrey Shi](https://www.humphreyshi.com).
 
 Text2Video-Zero enables zero-shot video generation using either:

docs/source/en/api/pipelines/unidiffuser.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # UniDiffuser
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 The UniDiffuser model was proposed in [One Transformer Fits All Distributions in Multi-Modal Diffusion at Scale](https://huggingface.co/papers/2303.06555) by Fan Bao, Shen Nie, Kaiwen Xue, Chongxuan Li, Shi Pu, Yaole Wang, Gang Yue, Yue Cao, Hang Su, Jun Zhu.
 
 The abstract from the paper is:

docs/source/en/api/pipelines/wuerstchen.md

@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
 
 # Würstchen
 
+<div class="flex flex-wrap space-x-1">
+  <img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
+</div>
+
 <img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
 
 [Wuerstchen: An Efficient Architecture for Large-Scale Text-to-Image Diffusion Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.

docs/source/en/api/utilities.md

@@ -41,3 +41,11 @@ Utility and helper functions for working with 🤗 Diffusers.
 ## randn_tensor
 
 [[autodoc]] utils.torch_utils.randn_tensor
+
+## apply_layerwise_casting
+
+[[autodoc]] hooks.layerwise_casting.apply_layerwise_casting
+
+## apply_group_offloading
+
+[[autodoc]] hooks.group_offloading.apply_group_offloading

docs/source/en/installation.md

@@ -23,32 +23,60 @@ You should install 🤗 Diffusers in a [virtual environment](https://docs.python
 If you're unfamiliar with Python virtual environments, take a look at this [guide](https://packaging.python.org/guides/installing-using-pip-and-virtual-environments/).
 A virtual environment makes it easier to manage different projects and avoid compatibility issues between dependencies.
 
-Start by creating a virtual environment in your project directory:
+Create a virtual environment with Python or [uv](https://docs.astral.sh/uv/) (refer to [Installation](https://docs.astral.sh/uv/getting-started/installation/) for installation instructions), a fast Rust-based Python package and project manager.
+
+<hfoptions id="install">
+<hfoption id="uv">
 
 ```bash
-python -m venv .env
+uv venv my-env
+source my-env/bin/activate
 ```
 
-Activate the virtual environment:
+</hfoption>
+<hfoption id="Python">
 
 ```bash
-source .env/bin/activate
+python -m venv my-env
+source my-env/bin/activate
 ```
 
-You should also install 🤗 Transformers because 🤗 Diffusers relies on its models:
+</hfoption>
+</hfoptions>
+
+You should also install 🤗 Transformers because 🤗 Diffusers relies on its models.
 
 
 <frameworkcontent>
 <pt>
-Note - PyTorch only supports Python 3.8 - 3.11 on Windows.
+
+PyTorch only supports Python 3.8 - 3.11 on Windows. Install Diffusers with uv.
+
+```bash
+uv install diffusers["torch"] transformers
+```
+
+You can also install Diffusers with pip.
+
 ```bash
 pip install diffusers["torch"] transformers
 ```
+
 </pt>
 <jax>
+
+Install Diffusers with uv.
+
+```bash
+uv pip install diffusers["flax"] transformers
+```
+
+You can also install Diffusers with pip.
+
 ```bash
 pip install diffusers["flax"] transformers
 ```
+
 </jax>
 </frameworkcontent>
 

docs/source/en/optimization/memory.md

@@ -158,6 +158,83 @@ In order to properly offload models after they're called, it is required to run
 
 </Tip>
 
+## Group offloading
+
+Group offloading is the middle ground between sequential and model offloading. It works by offloading groups of internal layers (either `torch.nn.ModuleList` or `torch.nn.Sequential`), which uses less memory than model-level offloading. It is also faster than sequential-level offloading because the number of device synchronizations is reduced.
+
+To enable group offloading, call the [`~ModelMixin.enable_group_offload`] method on the model if it is a Diffusers model implementation. For any other model implementation, use [`~hooks.group_offloading.apply_group_offloading`]:
+
+```python
+import torch
+from diffusers import CogVideoXPipeline
+from diffusers.hooks import apply_group_offloading
+from diffusers.utils import export_to_video
+
+# Load the pipeline
+onload_device = torch.device("cuda")
+offload_device = torch.device("cpu")
+pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
+
+# We can utilize the enable_group_offload method for Diffusers model implementations
+pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True)
+
+# For any other model implementations, the apply_group_offloading function can be used
+apply_group_offloading(pipe.text_encoder, onload_device=onload_device, offload_type="block_level", num_blocks_per_group=2)
+apply_group_offloading(pipe.vae, onload_device=onload_device, offload_type="leaf_level")
+
+prompt = (
+    "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
+    "The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
+    "pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
+    "casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
+    "The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
+    "atmosphere of this unique musical performance."
+)
+video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
+# This utilized about 14.79 GB. It can be further reduced by using tiling and using leaf_level offloading throughout the pipeline.
+print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
+export_to_video(video, "output.mp4", fps=8)
+```
+
+Group offloading (for CUDA devices with support for asynchronous data transfer streams) overlaps data transfer and computation to reduce the overall execution time compared to sequential offloading. This is enabled using layer prefetching with CUDA streams. The next layer to be executed is loaded onto the accelerator device while the current layer is being executed - this increases the memory requirements slightly. Group offloading also supports leaf-level offloading (equivalent to sequential CPU offloading) but can be made much faster when using streams.
+
+## FP8 layerwise weight-casting
+
+PyTorch supports `torch.float8_e4m3fn` and `torch.float8_e5m2` as weight storage dtypes, but they can't be used for computation in many different tensor operations due to unimplemented kernel support. However, you can use these dtypes to store model weights in fp8 precision and upcast them on-the-fly when the layers are used in the forward pass. This is known as layerwise weight-casting.
+
+Typically, inference on most models is done with `torch.float16` or `torch.bfloat16` weight/computation precision. Layerwise weight-casting cuts down the memory footprint of the model weights by approximately half.
+
+```python
+import torch
+from diffusers import CogVideoXPipeline, CogVideoXTransformer3DModel
+from diffusers.utils import export_to_video
+
+model_id = "THUDM/CogVideoX-5b"
+
+# Load the model in bfloat16 and enable layerwise casting
+transformer = CogVideoXTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
+transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
+
+# Load the pipeline
+pipe = CogVideoXPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=torch.bfloat16)
+pipe.to("cuda")
+
+prompt = (
+    "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
+    "The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
+    "pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
+    "casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
+    "The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
+    "atmosphere of this unique musical performance."
+)
+video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
+export_to_video(video, "output.mp4", fps=8)
+```
+
+In the above example, layerwise casting is enabled on the transformer component of the pipeline. By default, certain layers are skipped from the FP8 weight casting because it can lead to significant degradation of generation quality. The normalization and modulation related weight parameters are also skipped by default.
+
+However, you gain more control and flexibility by directly utilizing the [`~hooks.layerwise_casting.apply_layerwise_casting`] function instead of [`~ModelMixin.enable_layerwise_casting`].
+
 ## Channels-last memory format
 
 The channels-last memory format is an alternative way of ordering NCHW tensors in memory to preserve dimension ordering. Channels-last tensors are ordered in such a way that the channels become the densest dimension (storing images pixel-per-pixel). Since not all operators currently support the channels-last format, it may result in worst performance but you should still try and see if it works for your model.

docs/source/en/optimization/para_attn.md

@@ -0,0 +1,497 @@
+# ParaAttention
+
+<div class="flex justify-center">
+    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-performance.png">
+</div>
+<div class="flex justify-center">
+    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-performance.png">
+</div>
+
+
+Large image and video generation models, such as [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev) and [HunyuanVideo](https://huggingface.co/tencent/HunyuanVideo), can be an inference challenge for real-time applications and deployment because of their size.
+
+[ParaAttention](https://github.com/chengzeyi/ParaAttention) is a library that implements **context parallelism** and **first block cache**, and can be combined with other techniques (torch.compile, fp8 dynamic quantization), to accelerate inference.
+
+This guide will show you how to apply ParaAttention to FLUX.1-dev and HunyuanVideo on NVIDIA L20 GPUs.
+No optimizations are applied for our baseline benchmark, except for HunyuanVideo to avoid out-of-memory errors.
+
+Our baseline benchmark shows that FLUX.1-dev is able to generate a 1024x1024 resolution image in 28 steps in 26.36 seconds, and HunyuanVideo is able to generate 129 frames at 720p resolution in 30 steps in 3675.71 seconds.
+
+> [!TIP]
+> For even faster inference with context parallelism, try using NVIDIA A100 or H100 GPUs (if available) with NVLink support, especially when there is a large number of GPUs.
+
+## First Block Cache
+
+Caching the output of the transformers blocks in the model and reusing them in the next inference steps reduces the computation cost and makes inference faster.
+
+However, it is hard to decide when to reuse the cache to ensure quality generated images or videos. ParaAttention directly uses the **residual difference of the first transformer block output** to approximate the difference among model outputs. When the difference is small enough, the residual difference of previous inference steps is reused. In other words, the denoising step is skipped.
+
+This achieves a 2x speedup on FLUX.1-dev and HunyuanVideo inference with very good quality.
+
+<figure>
+    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/ada-cache.png" alt="Cache in Diffusion Transformer" />
+    <figcaption>How AdaCache works, First Block Cache is a variant of it</figcaption>
+</figure>
+
+<hfoptions id="first-block-cache">
+<hfoption id="FLUX-1.dev">
+
+To apply first block cache on FLUX.1-dev, call `apply_cache_on_pipe` as shown below. 0.08 is the default residual difference value for FLUX models.
+
+```python
+import time
+import torch
+from diffusers import FluxPipeline
+
+pipe = FluxPipeline.from_pretrained(
+    "black-forest-labs/FLUX.1-dev",
+    torch_dtype=torch.bfloat16,
+).to("cuda")
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(pipe, residual_diff_threshold=0.08)
+
+# Enable memory savings
+# pipe.enable_model_cpu_offload()
+# pipe.enable_sequential_cpu_offload()
+
+begin = time.time()
+image = pipe(
+    "A cat holding a sign that says hello world",
+    num_inference_steps=28,
+).images[0]
+end = time.time()
+print(f"Time: {end - begin:.2f}s")
+
+print("Saving image to flux.png")
+image.save("flux.png")
+```
+
+| Optimizations | Original | FBCache rdt=0.06 | FBCache rdt=0.08 | FBCache rdt=0.10 | FBCache rdt=0.12 |
+| - | - | - | - | - | - |
+| Preview | ![Original](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-original.png) | ![FBCache rdt=0.06](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.06.png) | ![FBCache rdt=0.08](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.08.png) | ![FBCache rdt=0.10](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.10.png) | ![FBCache rdt=0.12](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.12.png) |
+| Wall Time (s) | 26.36 | 21.83 | 17.01 | 16.00 | 13.78 |
+
+First Block Cache reduced the inference speed to 17.01 seconds compared to the baseline, or 1.55x faster, while maintaining nearly zero quality loss.
+
+</hfoption>
+<hfoption id="HunyuanVideo">
+
+To apply First Block Cache on HunyuanVideo, `apply_cache_on_pipe` as shown below. 0.06 is the default residual difference value for HunyuanVideo models.
+
+```python
+import time
+import torch
+from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
+from diffusers.utils import export_to_video
+
+model_id = "tencent/HunyuanVideo"
+transformer = HunyuanVideoTransformer3DModel.from_pretrained(
+    model_id,
+    subfolder="transformer",
+    torch_dtype=torch.bfloat16,
+    revision="refs/pr/18",
+)
+pipe = HunyuanVideoPipeline.from_pretrained(
+    model_id,
+    transformer=transformer,
+    torch_dtype=torch.float16,
+    revision="refs/pr/18",
+).to("cuda")
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(pipe, residual_diff_threshold=0.6)
+
+pipe.vae.enable_tiling()
+
+begin = time.time()
+output = pipe(
+    prompt="A cat walks on the grass, realistic",
+    height=720,
+    width=1280,
+    num_frames=129,
+    num_inference_steps=30,
+).frames[0]
+end = time.time()
+print(f"Time: {end - begin:.2f}s")
+
+print("Saving video to hunyuan_video.mp4")
+export_to_video(output, "hunyuan_video.mp4", fps=15)
+```
+
+<video controls>
+  <source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-original.mp4" type="video/mp4">
+  Your browser does not support the video tag.
+</video>
+
+<small> HunyuanVideo without FBCache </small>
+
+<video controls>
+  <source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-fbc.mp4" type="video/mp4">
+  Your browser does not support the video tag.
+</video>
+
+<small> HunyuanVideo with FBCache </small>
+
+First Block Cache reduced the inference speed to 2271.06 seconds compared to the baseline, or 1.62x faster, while maintaining nearly zero quality loss.
+
+</hfoption>
+</hfoptions>
+
+## fp8 quantization
+
+fp8 with dynamic quantization further speeds up inference and reduces memory usage. Both the activations and weights must be quantized in order to use the 8-bit [NVIDIA Tensor Cores](https://www.nvidia.com/en-us/data-center/tensor-cores/).
+
+Use `float8_weight_only` and `float8_dynamic_activation_float8_weight` to quantize the text encoder and transformer model.
+
+The default quantization method is per tensor quantization, but if your GPU supports row-wise quantization, you can also try it for better accuracy.
+
+Install [torchao](https://github.com/pytorch/ao/tree/main) with the command below.
+
+```bash
+pip3 install -U torch torchao
+```
+
+[torch.compile](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) with `mode="max-autotune-no-cudagraphs"` or `mode="max-autotune"` selects the best kernel for performance. Compilation can take a long time if it's the first time the model is called, but it is worth it once the model has been compiled.
+
+This example only quantizes the transformer model, but you can also quantize the text encoder to reduce memory usage even more.
+
+> [!TIP]
+> Dynamic quantization can significantly change the distribution of the model output, so you need to change the `residual_diff_threshold` to a larger value for it to take effect.
+
+<hfoptions id="fp8-quantization">
+<hfoption id="FLUX-1.dev">
+
+```python
+import time
+import torch
+from diffusers import FluxPipeline
+
+pipe = FluxPipeline.from_pretrained(
+    "black-forest-labs/FLUX.1-dev",
+    torch_dtype=torch.bfloat16,
+).to("cuda")
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(
+    pipe,
+    residual_diff_threshold=0.12,  # Use a larger value to make the cache take effect
+)
+
+from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
+
+quantize_(pipe.text_encoder, float8_weight_only())
+quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
+pipe.transformer = torch.compile(
+   pipe.transformer, mode="max-autotune-no-cudagraphs",
+)
+
+# Enable memory savings
+# pipe.enable_model_cpu_offload()
+# pipe.enable_sequential_cpu_offload()
+
+for i in range(2):
+    begin = time.time()
+    image = pipe(
+        "A cat holding a sign that says hello world",
+        num_inference_steps=28,
+    ).images[0]
+    end = time.time()
+    if i == 0:
+        print(f"Warm up time: {end - begin:.2f}s")
+    else:
+        print(f"Time: {end - begin:.2f}s")
+
+print("Saving image to flux.png")
+image.save("flux.png")
+```
+
+fp8 dynamic quantization and torch.compile reduced the inference speed to 7.56 seconds compared to the baseline, or 3.48x faster.
+
+</hfoption>
+<hfoption id="HunyuanVideo">
+
+```python
+import time
+import torch
+from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
+from diffusers.utils import export_to_video
+
+model_id = "tencent/HunyuanVideo"
+transformer = HunyuanVideoTransformer3DModel.from_pretrained(
+    model_id,
+    subfolder="transformer",
+    torch_dtype=torch.bfloat16,
+    revision="refs/pr/18",
+)
+pipe = HunyuanVideoPipeline.from_pretrained(
+    model_id,
+    transformer=transformer,
+    torch_dtype=torch.float16,
+    revision="refs/pr/18",
+).to("cuda")
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(pipe)
+
+from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
+
+quantize_(pipe.text_encoder, float8_weight_only())
+quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
+pipe.transformer = torch.compile(
+   pipe.transformer, mode="max-autotune-no-cudagraphs",
+)
+
+# Enable memory savings
+pipe.vae.enable_tiling()
+# pipe.enable_model_cpu_offload()
+# pipe.enable_sequential_cpu_offload()
+
+for i in range(2):
+    begin = time.time()
+    output = pipe(
+        prompt="A cat walks on the grass, realistic",
+        height=720,
+        width=1280,
+        num_frames=129,
+        num_inference_steps=1 if i == 0 else 30,
+    ).frames[0]
+    end = time.time()
+    if i == 0:
+        print(f"Warm up time: {end - begin:.2f}s")
+    else:
+        print(f"Time: {end - begin:.2f}s")
+
+print("Saving video to hunyuan_video.mp4")
+export_to_video(output, "hunyuan_video.mp4", fps=15)
+```
+
+A NVIDIA L20 GPU only has 48GB memory and could face out-of-memory (OOM) errors after compilation and if `enable_model_cpu_offload` isn't called because HunyuanVideo has very large activation tensors when running with high resolution and large number of frames. For GPUs with less than 80GB of memory, you can try reducing the resolution and number of frames to avoid OOM errors.
+
+Large video generation models are usually bottlenecked by the attention computations rather than the fully connected layers. These models don't significantly benefit from quantization and torch.compile.
+
+</hfoption>
+</hfoptions>
+
+## Context Parallelism
+
+Context Parallelism parallelizes inference and scales with multiple GPUs. The ParaAttention compositional design allows you to combine Context Parallelism with First Block Cache and dynamic quantization.
+
+> [!TIP]
+> Refer to the [ParaAttention](https://github.com/chengzeyi/ParaAttention/tree/main) repository for detailed instructions and examples of how to scale inference with multiple GPUs.
+
+If the inference process needs to be persistent and serviceable, it is suggested to use [torch.multiprocessing](https://pytorch.org/docs/stable/multiprocessing.html) to write your own inference processor. This can eliminate the overhead of launching the process and loading and recompiling the model.
+
+<hfoptions id="context-parallelism">
+<hfoption id="FLUX-1.dev">
+
+The code sample below combines First Block Cache, fp8 dynamic quantization, torch.compile, and Context Parallelism for the fastest inference speed.
+
+```python
+import time
+import torch
+import torch.distributed as dist
+from diffusers import FluxPipeline
+
+dist.init_process_group()
+
+torch.cuda.set_device(dist.get_rank())
+
+pipe = FluxPipeline.from_pretrained(
+    "black-forest-labs/FLUX.1-dev",
+    torch_dtype=torch.bfloat16,
+).to("cuda")
+
+from para_attn.context_parallel import init_context_parallel_mesh
+from para_attn.context_parallel.diffusers_adapters import parallelize_pipe
+from para_attn.parallel_vae.diffusers_adapters import parallelize_vae
+
+mesh = init_context_parallel_mesh(
+    pipe.device.type,
+    max_ring_dim_size=2,
+)
+parallelize_pipe(
+    pipe,
+    mesh=mesh,
+)
+parallelize_vae(pipe.vae, mesh=mesh._flatten())
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(
+    pipe,
+    residual_diff_threshold=0.12,  # Use a larger value to make the cache take effect
+)
+
+from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
+
+quantize_(pipe.text_encoder, float8_weight_only())
+quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
+torch._inductor.config.reorder_for_compute_comm_overlap = True
+pipe.transformer = torch.compile(
+   pipe.transformer, mode="max-autotune-no-cudagraphs",
+)
+
+# Enable memory savings
+# pipe.enable_model_cpu_offload(gpu_id=dist.get_rank())
+# pipe.enable_sequential_cpu_offload(gpu_id=dist.get_rank())
+
+for i in range(2):
+    begin = time.time()
+    image = pipe(
+        "A cat holding a sign that says hello world",
+        num_inference_steps=28,
+        output_type="pil" if dist.get_rank() == 0 else "pt",
+    ).images[0]
+    end = time.time()
+    if dist.get_rank() == 0:
+        if i == 0:
+            print(f"Warm up time: {end - begin:.2f}s")
+        else:
+            print(f"Time: {end - begin:.2f}s")
+
+if dist.get_rank() == 0:
+    print("Saving image to flux.png")
+    image.save("flux.png")
+
+dist.destroy_process_group()
+```
+
+Save to `run_flux.py` and launch it with [torchrun](https://pytorch.org/docs/stable/elastic/run.html).
+
+```bash
+# Use --nproc_per_node to specify the number of GPUs
+torchrun --nproc_per_node=2 run_flux.py
+```
+
+Inference speed is reduced to 8.20 seconds compared to the baseline, or 3.21x faster, with 2 NVIDIA L20 GPUs. On 4 L20s, inference speed is 3.90 seconds, or 6.75x faster.
+
+</hfoption>
+<hfoption id="HunyuanVideo">
+
+The code sample below combines First Block Cache and Context Parallelism for the fastest inference speed.
+
+```python
+import time
+import torch
+import torch.distributed as dist
+from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
+from diffusers.utils import export_to_video
+
+dist.init_process_group()
+
+torch.cuda.set_device(dist.get_rank())
+
+model_id = "tencent/HunyuanVideo"
+transformer = HunyuanVideoTransformer3DModel.from_pretrained(
+    model_id,
+    subfolder="transformer",
+    torch_dtype=torch.bfloat16,
+    revision="refs/pr/18",
+)
+pipe = HunyuanVideoPipeline.from_pretrained(
+    model_id,
+    transformer=transformer,
+    torch_dtype=torch.float16,
+    revision="refs/pr/18",
+).to("cuda")
+
+from para_attn.context_parallel import init_context_parallel_mesh
+from para_attn.context_parallel.diffusers_adapters import parallelize_pipe
+from para_attn.parallel_vae.diffusers_adapters import parallelize_vae
+
+mesh = init_context_parallel_mesh(
+    pipe.device.type,
+)
+parallelize_pipe(
+    pipe,
+    mesh=mesh,
+)
+parallelize_vae(pipe.vae, mesh=mesh._flatten())
+
+from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
+
+apply_cache_on_pipe(pipe)
+
+# from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
+#
+# torch._inductor.config.reorder_for_compute_comm_overlap = True
+#
+# quantize_(pipe.text_encoder, float8_weight_only())
+# quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
+# pipe.transformer = torch.compile(
+#    pipe.transformer, mode="max-autotune-no-cudagraphs",
+# )
+
+# Enable memory savings
+pipe.vae.enable_tiling()
+# pipe.enable_model_cpu_offload(gpu_id=dist.get_rank())
+# pipe.enable_sequential_cpu_offload(gpu_id=dist.get_rank())
+
+for i in range(2):
+    begin = time.time()
+    output = pipe(
+        prompt="A cat walks on the grass, realistic",
+        height=720,
+        width=1280,
+        num_frames=129,
+        num_inference_steps=1 if i == 0 else 30,
+        output_type="pil" if dist.get_rank() == 0 else "pt",
+    ).frames[0]
+    end = time.time()
+    if dist.get_rank() == 0:
+        if i == 0:
+            print(f"Warm up time: {end - begin:.2f}s")
+        else:
+            print(f"Time: {end - begin:.2f}s")
+
+if dist.get_rank() == 0:
+    print("Saving video to hunyuan_video.mp4")
+    export_to_video(output, "hunyuan_video.mp4", fps=15)
+
+dist.destroy_process_group()
+```
+
+Save to `run_hunyuan_video.py` and launch it with [torchrun](https://pytorch.org/docs/stable/elastic/run.html).
+
+```bash
+# Use --nproc_per_node to specify the number of GPUs
+torchrun --nproc_per_node=8 run_hunyuan_video.py
+```
+
+Inference speed is reduced to 649.23 seconds compared to the baseline, or 5.66x faster, with 8 NVIDIA L20 GPUs.
+
+</hfoption>
+</hfoptions>
+
+## Benchmarks
+
+<hfoptions id="conclusion">
+<hfoption id="FLUX-1.dev">
+
+| GPU Type | Number of GPUs | Optimizations | Wall Time (s) | Speedup |
+| - | - | - | - | - |
+| NVIDIA L20 | 1 | Baseline | 26.36 | 1.00x |
+| NVIDIA L20 | 1 | FBCache (rdt=0.08) | 17.01 | 1.55x |
+| NVIDIA L20 | 1 | FP8 DQ | 13.40 | 1.96x |
+| NVIDIA L20 | 1 | FBCache (rdt=0.12) + FP8 DQ | 7.56 | 3.48x |
+| NVIDIA L20 | 2 | FBCache (rdt=0.12) + FP8 DQ + CP | 4.92 | 5.35x |
+| NVIDIA L20 | 4 | FBCache (rdt=0.12) + FP8 DQ + CP | 3.90 | 6.75x |
+
+</hfoption>
+<hfoption id="HunyuanVideo">
+
+| GPU Type | Number of GPUs | Optimizations | Wall Time (s) | Speedup |
+| - | - | - | - | - |
+| NVIDIA L20 | 1 | Baseline | 3675.71 | 1.00x |
+| NVIDIA L20 | 1 | FBCache | 2271.06 | 1.62x |
+| NVIDIA L20 | 2 | FBCache + CP | 1132.90 | 3.24x |
+| NVIDIA L20 | 4 | FBCache + CP | 718.15 | 5.12x |
+| NVIDIA L20 | 8 | FBCache + CP | 649.23 | 5.66x |
+
+</hfoption>
+</hfoptions>

docs/source/en/training/custom_diffusion.md

@@ -339,7 +339,10 @@ import torch
 from huggingface_hub.repocard import RepoCard
 from diffusers import DiffusionPipeline
 
-pipeline = DiffusionPipeline.from_pretrained("sayakpaul/custom-diffusion-cat-wooden-pot", torch_dtype=torch.float16).to("cuda")
+pipeline = DiffusionPipeline.from_pretrained(
+    "CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16,
+).to("cuda")
+model_id = "sayakpaul/custom-diffusion-cat-wooden-pot"
 pipeline.unet.load_attn_procs(model_id, weight_name="pytorch_custom_diffusion_weights.bin")
 pipeline.load_textual_inversion(model_id, weight_name="<new1>.bin")
 pipeline.load_textual_inversion(model_id, weight_name="<new2>.bin")

docs/source/en/tutorials/using_peft_for_inference.md

@@ -221,3 +221,7 @@ pipe.delete_adapters("toy")
 pipe.get_active_adapters()
 ["pixel"]
 ```
+
+## PeftInputAutocastDisableHook
+
+[[autodoc]] hooks.layerwise_casting.PeftInputAutocastDisableHook

docs/source/en/using-diffusers/consisid.md

@@ -0,0 +1,96 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License.
+-->
+# ConsisID
+
+[ConsisID](https://github.com/PKU-YuanGroup/ConsisID) is an identity-preserving text-to-video generation model that keeps the face consistent in the generated video by frequency decomposition. The main features of ConsisID are:
+
+- Frequency decomposition: The characteristics of the DiT architecture are analyzed from the frequency domain perspective, and based on these characteristics, a reasonable control information injection method is designed.
+- Consistency training strategy: A coarse-to-fine training strategy, dynamic masking loss, and dynamic cross-face loss further enhance the model's generalization ability and identity preservation performance.
+- Inference without finetuning: Previous methods required case-by-case finetuning of the input ID before inference, leading to significant time and computational costs. In contrast, ConsisID is tuning-free.
+
+This guide will walk you through using ConsisID for use cases.
+
+## Load Model Checkpoints
+
+Model weights may be stored in separate subfolders on the Hub or locally, in which case, you should use the [`~DiffusionPipeline.from_pretrained`] method.
+
+```python
+# !pip install consisid_eva_clip insightface facexlib
+import torch
+from diffusers import ConsisIDPipeline
+from diffusers.pipelines.consisid.consisid_utils import prepare_face_models, process_face_embeddings_infer
+from huggingface_hub import snapshot_download
+
+# Download ckpts
+snapshot_download(repo_id="BestWishYsh/ConsisID-preview", local_dir="BestWishYsh/ConsisID-preview")
+
+# Load face helper model to preprocess input face image
+face_helper_1, face_helper_2, face_clip_model, face_main_model, eva_transform_mean, eva_transform_std = prepare_face_models("BestWishYsh/ConsisID-preview", device="cuda", dtype=torch.bfloat16)
+
+# Load consisid base model
+pipe = ConsisIDPipeline.from_pretrained("BestWishYsh/ConsisID-preview", torch_dtype=torch.bfloat16)
+pipe.to("cuda")
+```
+
+## Identity-Preserving Text-to-Video
+
+For identity-preserving text-to-video, pass a text prompt and an image contain clear face (e.g., preferably half-body or full-body). By default, ConsisID generates a 720x480 video for the best results.
+
+```python
+from diffusers.utils import export_to_video
+
+prompt = "The video captures a boy walking along a city street, filmed in black and white on a classic 35mm camera. His expression is thoughtful, his brow slightly furrowed as if he's lost in contemplation. The film grain adds a textured, timeless quality to the image, evoking a sense of nostalgia. Around him, the cityscape is filled with vintage buildings, cobblestone sidewalks, and softly blurred figures passing by, their outlines faint and indistinct. Streetlights cast a gentle glow, while shadows play across the boy's path, adding depth to the scene. The lighting highlights the boy's subtle smile, hinting at a fleeting moment of curiosity. The overall cinematic atmosphere, complete with classic film still aesthetics and dramatic contrasts, gives the scene an evocative and introspective feel."
+image = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_input.png?download=true"
+
+id_cond, id_vit_hidden, image, face_kps = process_face_embeddings_infer(face_helper_1, face_clip_model, face_helper_2, eva_transform_mean, eva_transform_std, face_main_model, "cuda", torch.bfloat16, image, is_align_face=True)
+
+video = pipe(image=image, prompt=prompt, num_inference_steps=50, guidance_scale=6.0, use_dynamic_cfg=False, id_vit_hidden=id_vit_hidden, id_cond=id_cond, kps_cond=face_kps, generator=torch.Generator("cuda").manual_seed(42))
+export_to_video(video.frames[0], "output.mp4", fps=8)
+```
+<table>
+  <tr>
+    <th style="text-align: center;">Face Image</th>
+    <th style="text-align: center;">Video</th>
+    <th style="text-align: center;">Description</th
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_0.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_0.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video, in a beautifully crafted animated style, features a confident woman riding a horse through a lush forest clearing. Her expression is focused yet serene as she adjusts her wide-brimmed hat with a practiced hand. She wears a flowy bohemian dress, which moves gracefully with the rhythm of the horse, the fabric flowing fluidly in the animated motion. The dappled sunlight filters through the trees, casting soft, painterly patterns on the forest floor. Her posture is poised, showing both control and elegance as she guides the horse with ease. The animation's gentle, fluid style adds a dreamlike quality to the scene, with the woman’s calm demeanor and the peaceful surroundings evoking a sense of freedom and harmony.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_1.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_1.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video, in a captivating animated style, shows a woman standing in the center of a snowy forest, her eyes narrowed in concentration as she extends her hand forward. She is dressed in a deep blue cloak, her breath visible in the cold air, which is rendered with soft, ethereal strokes. A faint smile plays on her lips as she summons a wisp of ice magic, watching with focus as the surrounding trees and ground begin to shimmer and freeze, covered in delicate ice crystals. The animation’s fluid motion brings the magic to life, with the frost spreading outward in intricate, sparkling patterns. The environment is painted with soft, watercolor-like hues, enhancing the magical, dreamlike atmosphere. The overall mood is serene yet powerful, with the quiet winter air amplifying the delicate beauty of the frozen scene.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_2.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_2.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The animation features a whimsical portrait of a balloon seller standing in a gentle breeze, captured with soft, hazy brushstrokes that evoke the feel of a serene spring day. His face is framed by a gentle smile, his eyes squinting slightly against the sun, while a few wisps of hair flutter in the wind. He is dressed in a light, pastel-colored shirt, and the balloons around him sway with the wind, adding a sense of playfulness to the scene. The background blurs softly, with hints of a vibrant market or park, enhancing the light-hearted, yet tender mood of the moment.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_3.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_3.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video captures a boy walking along a city street, filmed in black and white on a classic 35mm camera. His expression is thoughtful, his brow slightly furrowed as if he's lost in contemplation. The film grain adds a textured, timeless quality to the image, evoking a sense of nostalgia. Around him, the cityscape is filled with vintage buildings, cobblestone sidewalks, and softly blurred figures passing by, their outlines faint and indistinct. Streetlights cast a gentle glow, while shadows play across the boy's path, adding depth to the scene. The lighting highlights the boy's subtle smile, hinting at a fleeting moment of curiosity. The overall cinematic atmosphere, complete with classic film still aesthetics and dramatic contrasts, gives the scene an evocative and introspective feel.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_4.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_4.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video features a baby wearing a bright superhero cape, standing confidently with arms raised in a powerful pose. The baby has a determined look on their face, with eyes wide and lips pursed in concentration, as if ready to take on a challenge. The setting appears playful, with colorful toys scattered around and a soft rug underfoot, while sunlight streams through a nearby window, highlighting the fluttering cape and adding to the impression of heroism. The overall atmosphere is lighthearted and fun, with the baby's expressions capturing a mix of innocence and an adorable attempt at bravery, as if truly ready to save the day.</td>
+  </tr>
+</table>
+
+## Resources
+
+Learn more about ConsisID with the following resources.
+- A [video](https://www.youtube.com/watch?v=PhlgC-bI5SQ) demonstrating ConsisID's main features.
+- The research paper, [Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://hf.co/papers/2411.17440) for more details.

docs/source/en/using-diffusers/img2img.md

@@ -461,12 +461,12 @@ Chain it to an upscaler pipeline to increase the image resolution:
 from diffusers import StableDiffusionLatentUpscalePipeline
 
 upscaler = StableDiffusionLatentUpscalePipeline.from_pretrained(
-    "stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
+    "stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16, use_safetensors=True
 )
 upscaler.enable_model_cpu_offload()
 upscaler.enable_xformers_memory_efficient_attention()
 
-image_2 = upscaler(prompt, image=image_1, output_type="latent").images[0]
+image_2 = upscaler(prompt, image=image_1).images[0]
 ```
 
 Finally, chain it to a super-resolution pipeline to further enhance the resolution:

docs/source/en/using-diffusers/omnigen.md

@@ -0,0 +1,317 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License.
+-->
+# OmniGen
+
+OmniGen is an image generation model. Unlike existing text-to-image models, OmniGen is a single model designed to handle a variety of tasks (e.g., text-to-image, image editing, controllable generation). It has the following features:
+- Minimalist model architecture, consisting of only a VAE and a transformer module, for joint modeling of text and images.
+- Support for multimodal inputs. It can process any text-image mixed data as instructions for image generation, rather than relying solely on text.
+
+For more information, please refer to the [paper](https://arxiv.org/pdf/2409.11340).
+This guide will walk you through using OmniGen for various tasks and use cases.
+
+## Load model checkpoints
+
+Model weights may be stored in separate subfolders on the Hub or locally, in which case, you should use the [`~DiffusionPipeline.from_pretrained`] method.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+
+pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
+```
+
+## Text-to-image
+
+For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image. 
+You can try setting the `height` and `width` parameters to generate images with different size.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
+image = pipe(
+    prompt=prompt,
+    height=1024,
+    width=1024,
+    guidance_scale=3,
+    generator=torch.Generator(device="cpu").manual_seed(111),
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex justify-center">
+    <img src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png" alt="generated image"/>
+</div>
+
+## Image edit
+
+OmniGen supports multimodal inputs. 
+When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image. 
+It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+from diffusers.utils import load_image 
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(222)
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex flex-row gap-4">
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">edited image</figcaption>
+  </div>
+</div>
+
+OmniGen has some interesting features, such as visual reasoning, as shown in the example below.
+
+```python
+prompt="If the woman is thirsty, what should she take? Find it in the image and highlight it in blue. <img><|image_1|></img>"
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(0)
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex justify-center">
+    <img src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/reasoning.png" alt="generated image"/>
+</div>
+
+## Controllable generation
+
+OmniGen can handle several classic computer vision tasks. As shown below, OmniGen can detect human skeletons in input images, which can be used as control conditions to generate new images.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+from diffusers.utils import load_image 
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt="Detect the skeleton of human in this image: <img><|image_1|></img>"
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
+image1 = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(333)
+).images[0]
+image1.save("image1.png")
+
+prompt="Generate a new photo using the following picture and text as conditions: <img><|image_1|></img>\n A young boy is sitting on a sofa in the library, holding a book. His hair is neatly combed, and a faint smile plays on his lips, with a few freckles scattered across his cheeks. The library is quiet, with rows of shelves filled with books stretching out behind him."
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal.png")]
+image2 = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(333)
+).images[0]
+image2.save("image2.png")
+```
+
+<div class="flex flex-row gap-4">
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">detected skeleton</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal2img.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">skeleton to image</figcaption>
+  </div>
+</div>
+
+
+OmniGen can also directly use relevant information from input images to generate new images.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+from diffusers.utils import load_image 
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt="Following the pose of this image <img><|image_1|></img>, generate a new photo: A young boy is sitting on a sofa in the library, holding a book. His hair is neatly combed, and a faint smile plays on his lips, with a few freckles scattered across his cheeks. The library is quiet, with rows of shelves filled with books stretching out behind him."
+input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    guidance_scale=2, 
+    img_guidance_scale=1.6,
+    use_input_image_size_as_output=True,
+    generator=torch.Generator(device="cpu").manual_seed(0)
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex flex-row gap-4">
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/same_pose.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
+  </div>
+</div>
+
+## ID and object preserving
+
+OmniGen can generate multiple images based on the people and objects in the input image and supports inputting multiple images simultaneously. 
+Additionally, OmniGen can extract desired objects from an image containing multiple objects based on instructions.
+
+```python
+import torch
+from diffusers import OmniGenPipeline
+from diffusers.utils import load_image 
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt="A man and a woman are sitting at a classroom desk. The man is the man with yellow hair in <img><|image_1|></img>. The woman is the woman on the left of <img><|image_2|></img>"
+input_image_1 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/3.png")
+input_image_2 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/4.png")
+input_images=[input_image_1, input_image_2]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    height=1024,
+    width=1024,
+    guidance_scale=2.5, 
+    img_guidance_scale=1.6,
+    generator=torch.Generator(device="cpu").manual_seed(666)
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex flex-row gap-4">
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/3.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">input_image_1</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/4.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">input_image_2</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/id2.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
+  </div>
+</div>
+
+```py
+import torch
+from diffusers import OmniGenPipeline
+from diffusers.utils import load_image 
+
+pipe = OmniGenPipeline.from_pretrained(
+    "Shitao/OmniGen-v1-diffusers",
+    torch_dtype=torch.bfloat16
+)
+pipe.to("cuda")
+
+prompt="A woman is walking down the street, wearing a white long-sleeve blouse with lace details on the sleeves, paired with a blue pleated skirt. The woman is <img><|image_1|></img>. The long-sleeve blouse and a pleated skirt are <img><|image_2|></img>."
+input_image_1 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/emma.jpeg")
+input_image_2 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/dress.jpg")
+input_images=[input_image_1, input_image_2]
+image = pipe(
+    prompt=prompt, 
+    input_images=input_images, 
+    height=1024,
+    width=1024,
+    guidance_scale=2.5, 
+    img_guidance_scale=1.6,
+    generator=torch.Generator(device="cpu").manual_seed(666)
+).images[0]
+image.save("output.png")
+```
+
+<div class="flex flex-row gap-4">
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/emma.jpeg"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">person image</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/dress.jpg"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">clothe image</figcaption>
+  </div>
+  <div class="flex-1">
+    <img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/tryon.png"/>
+    <figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
+  </div>
+</div>
+
+## Optimization when using multiple images 
+
+For text-to-image task, OmniGen requires minimal memory and time costs (9GB memory and 31s for a 1024x1024 image on A800 GPU). 
+However, when using input images, the computational cost increases. 
+
+Here are some guidelines to help you reduce computational costs when using multiple images. The experiments are conducted on an A800 GPU with two input images.
+
+Like other pipelines, you can reduce memory usage by offloading the model: `pipe.enable_model_cpu_offload()` or `pipe.enable_sequential_cpu_offload() `. 
+In OmniGen, you can also decrease computational overhead by reducing the `max_input_image_size`. 
+The memory consumption for different image sizes is shown in the table below:
+
+| Method                    | Memory Usage |
+|---------------------------|--------------|
+| max_input_image_size=1024 | 40GB         |
+| max_input_image_size=512  | 17GB         |
+| max_input_image_size=256  | 14GB         |
+

docs/source/en/using-diffusers/other-formats.md

@@ -240,6 +240,46 @@ Benefits of using a single-file layout include:
 1. Easy compatibility with diffusion interfaces such as [ComfyUI](https://github.com/comfyanonymous/ComfyUI) or [Automatic1111](https://github.com/AUTOMATIC1111/stable-diffusion-webui) which commonly use a single-file layout.
 2. Easier to manage (download and share) a single file.
 
+### DDUF
+
+> [!WARNING]
+> DDUF is an experimental file format and APIs related to it can change in the future.
+
+DDUF (**D**DUF **D**iffusion **U**nified **F**ormat) is a file format designed to make storing, distributing, and using diffusion models much easier. Built on the ZIP file format, DDUF offers a standardized, efficient, and flexible way to package all parts of a diffusion model into a single, easy-to-manage file. It provides a balance between Diffusers multi-folder format and the widely popular single-file format.
+
+Learn more details about DDUF on the Hugging Face Hub [documentation](https://huggingface.co/docs/hub/dduf).
+
+Pass a checkpoint to the `dduf_file` parameter to load it in [`DiffusionPipeline`].
+
+```py
+from diffusers import DiffusionPipeline
+import torch
+
+pipe = DiffusionPipeline.from_pretrained(
+    "DDUF/FLUX.1-dev-DDUF", dduf_file="FLUX.1-dev.dduf", torch_dtype=torch.bfloat16
+).to("cuda")
+image = pipe(
+    "photo a cat holding a sign that says Diffusers", num_inference_steps=50, guidance_scale=3.5
+).images[0]
+image.save("cat.png")
+```
+
+To save a pipeline as a `.dduf` checkpoint, use the [`~huggingface_hub.export_folder_as_dduf`] utility, which takes care of all the necessary file-level validations.
+
+```py
+from huggingface_hub import export_folder_as_dduf
+from diffusers import DiffusionPipeline
+import torch 
+
+pipe = DiffusionPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16)
+
+save_folder = "flux-dev"
+pipe.save_pretrained("flux-dev")
+export_folder_as_dduf("flux-dev.dduf", folder_path=save_folder)
+
+> [!TIP]
+> Packaging and loading quantized checkpoints in the DDUF format is supported as long as they respect the multi-folder structure.
+
 ## Convert layout and files
 
 Diffusers provides many scripts and methods to convert storage layouts and file formats to enable broader support across the diffusion ecosystem.

docs/source/en/using-diffusers/write_own_pipeline.md

@@ -106,7 +106,7 @@ Let's try it out!
 
 ## Deconstruct the Stable Diffusion pipeline
 
-Stable Diffusion is a text-to-image *latent diffusion* model. It is called a latent diffusion model because it works with a lower-dimensional representation of the image instead of the actual pixel space, which makes it more memory efficient. The encoder compresses the image into a smaller representation, and a decoder to convert the compressed representation back into an image. For text-to-image models, you'll need a tokenizer and an encoder to generate text embeddings. From the previous example, you already know you need a UNet model and a scheduler.
+Stable Diffusion is a text-to-image *latent diffusion* model. It is called a latent diffusion model because it works with a lower-dimensional representation of the image instead of the actual pixel space, which makes it more memory efficient. The encoder compresses the image into a smaller representation, and a decoder converts the compressed representation back into an image. For text-to-image models, you'll need a tokenizer and an encoder to generate text embeddings. From the previous example, you already know you need a UNet model and a scheduler.
 
 As you can see, this is already more complex than the DDPM pipeline which only contains a UNet model. The Stable Diffusion model has three separate pretrained models.
 

docs/source/zh/_toctree.yml

@@ -5,6 +5,8 @@
     title: 快速入门
   - local: stable_diffusion
     title: 有效和高效的扩散
+  - local: consisid 
+    title: 身份保持的文本到视频生成
   - local: installation
     title: 安装
   title: 开始

docs/source/zh/consisid.md

@@ -0,0 +1,100 @@
+<!--Copyright 2024 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License.
+-->
+# ConsisID
+
+[ConsisID](https://github.com/PKU-YuanGroup/ConsisID)是一种身份保持的文本到视频生成模型，其通过频率分解在生成的视频中保持面部一致性。它具有以下特点：
+
+- 基于频率分解：将人物ID特征解耦为高频和低频部分，从频域的角度分析DIT架构的特性，并且基于此特性设计合理的控制信息注入方式。
+
+- 一致性训练策略：我们提出粗到细训练策略、动态掩码损失、动态跨脸损失，进一步提高了模型的泛化能力和身份保持效果。
+
+
+- 推理无需微调：之前的方法在推理前，需要对输入id进行case-by-case微调，时间和算力开销较大，而我们的方法是tuning-free的。
+
+
+本指南将指导您使用 ConsisID 生成身份保持的视频。
+
+## Load Model Checkpoints
+模型权重可以存储在Hub上或本地的单独子文件夹中，在这种情况下，您应该使用 [`~DiffusionPipeline.from_pretrained`] 方法。
+
+
+```python
+# !pip install consisid_eva_clip insightface facexlib
+import torch
+from diffusers import ConsisIDPipeline
+from diffusers.pipelines.consisid.consisid_utils import prepare_face_models, process_face_embeddings_infer
+from huggingface_hub import snapshot_download
+
+# Download ckpts
+snapshot_download(repo_id="BestWishYsh/ConsisID-preview", local_dir="BestWishYsh/ConsisID-preview")
+
+# Load face helper model to preprocess input face image
+face_helper_1, face_helper_2, face_clip_model, face_main_model, eva_transform_mean, eva_transform_std = prepare_face_models("BestWishYsh/ConsisID-preview", device="cuda", dtype=torch.bfloat16)
+
+# Load consisid base model
+pipe = ConsisIDPipeline.from_pretrained("BestWishYsh/ConsisID-preview", torch_dtype=torch.bfloat16)
+pipe.to("cuda")
+```
+
+## Identity-Preserving Text-to-Video
+对于身份保持的文本到视频生成，需要输入文本提示和包含清晰面部（例如，最好是半身或全身）的图像。默认情况下，ConsisID 会生成 720x480 的视频以获得最佳效果。
+
+```python
+from diffusers.utils import export_to_video
+
+prompt = "The video captures a boy walking along a city street, filmed in black and white on a classic 35mm camera. His expression is thoughtful, his brow slightly furrowed as if he's lost in contemplation. The film grain adds a textured, timeless quality to the image, evoking a sense of nostalgia. Around him, the cityscape is filled with vintage buildings, cobblestone sidewalks, and softly blurred figures passing by, their outlines faint and indistinct. Streetlights cast a gentle glow, while shadows play across the boy's path, adding depth to the scene. The lighting highlights the boy's subtle smile, hinting at a fleeting moment of curiosity. The overall cinematic atmosphere, complete with classic film still aesthetics and dramatic contrasts, gives the scene an evocative and introspective feel."
+image = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_input.png?download=true"
+
+id_cond, id_vit_hidden, image, face_kps = process_face_embeddings_infer(face_helper_1, face_clip_model, face_helper_2, eva_transform_mean, eva_transform_std, face_main_model, "cuda", torch.bfloat16, image, is_align_face=True)
+
+video = pipe(image=image, prompt=prompt, num_inference_steps=50, guidance_scale=6.0, use_dynamic_cfg=False, id_vit_hidden=id_vit_hidden, id_cond=id_cond, kps_cond=face_kps, generator=torch.Generator("cuda").manual_seed(42))
+export_to_video(video.frames[0], "output.mp4", fps=8)
+```
+<table>
+  <tr>
+    <th style="text-align: center;">Face Image</th>
+    <th style="text-align: center;">Video</th>
+    <th style="text-align: center;">Description</th
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_0.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_0.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video, in a beautifully crafted animated style, features a confident woman riding a horse through a lush forest clearing. Her expression is focused yet serene as she adjusts her wide-brimmed hat with a practiced hand. She wears a flowy bohemian dress, which moves gracefully with the rhythm of the horse, the fabric flowing fluidly in the animated motion. The dappled sunlight filters through the trees, casting soft, painterly patterns on the forest floor. Her posture is poised, showing both control and elegance as she guides the horse with ease. The animation's gentle, fluid style adds a dreamlike quality to the scene, with the woman’s calm demeanor and the peaceful surroundings evoking a sense of freedom and harmony.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_1.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_1.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video, in a captivating animated style, shows a woman standing in the center of a snowy forest, her eyes narrowed in concentration as she extends her hand forward. She is dressed in a deep blue cloak, her breath visible in the cold air, which is rendered with soft, ethereal strokes. A faint smile plays on her lips as she summons a wisp of ice magic, watching with focus as the surrounding trees and ground begin to shimmer and freeze, covered in delicate ice crystals. The animation’s fluid motion brings the magic to life, with the frost spreading outward in intricate, sparkling patterns. The environment is painted with soft, watercolor-like hues, enhancing the magical, dreamlike atmosphere. The overall mood is serene yet powerful, with the quiet winter air amplifying the delicate beauty of the frozen scene.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_2.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_2.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The animation features a whimsical portrait of a balloon seller standing in a gentle breeze, captured with soft, hazy brushstrokes that evoke the feel of a serene spring day. His face is framed by a gentle smile, his eyes squinting slightly against the sun, while a few wisps of hair flutter in the wind. He is dressed in a light, pastel-colored shirt, and the balloons around him sway with the wind, adding a sense of playfulness to the scene. The background blurs softly, with hints of a vibrant market or park, enhancing the light-hearted, yet tender mood of the moment.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_3.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_3.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video captures a boy walking along a city street, filmed in black and white on a classic 35mm camera. His expression is thoughtful, his brow slightly furrowed as if he's lost in contemplation. The film grain adds a textured, timeless quality to the image, evoking a sense of nostalgia. Around him, the cityscape is filled with vintage buildings, cobblestone sidewalks, and softly blurred figures passing by, their outlines faint and indistinct. Streetlights cast a gentle glow, while shadows play across the boy's path, adding depth to the scene. The lighting highlights the boy's subtle smile, hinting at a fleeting moment of curiosity. The overall cinematic atmosphere, complete with classic film still aesthetics and dramatic contrasts, gives the scene an evocative and introspective feel.</td>
+  </tr>
+  <tr>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_image_4.png?download=true" style="height: auto; width: 600px;"></td>
+    <td><img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/consisid/consisid_output_4.gif?download=true" style="height: auto; width: 2000px;"></td>
+    <td>The video features a baby wearing a bright superhero cape, standing confidently with arms raised in a powerful pose. The baby has a determined look on their face, with eyes wide and lips pursed in concentration, as if ready to take on a challenge. The setting appears playful, with colorful toys scattered around and a soft rug underfoot, while sunlight streams through a nearby window, highlighting the fluttering cape and adding to the impression of heroism. The overall atmosphere is lighthearted and fun, with the baby's expressions capturing a mix of innocence and an adorable attempt at bravery, as if truly ready to save the day.</td>
+  </tr>
+</table>
+
+## Resources
+
+通过以下资源了解有关 ConsisID 的更多信息：
+
+- 一段 [视频](https://www.youtube.com/watch?v=PhlgC-bI5SQ) 演示了 ConsisID 的主要功能；
+- 有关更多详细信息，请参阅研究论文 [Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://hf.co/papers/2411.17440)。

examples/README.md

@@ -40,9 +40,9 @@ Training examples show how to pretrain or fine-tune diffusion models for a varie
 | [**Text-to-Image fine-tuning**](./text_to_image) | ✅ | ✅ |
 | [**Textual Inversion**](./textual_inversion) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_textual_inversion_training.ipynb)
 | [**Dreambooth**](./dreambooth) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_dreambooth_training.ipynb)
-| [**ControlNet**](./controlnet) | ✅ | ✅ | -
-| [**InstructPix2Pix**](./instruct_pix2pix) | ✅ | ✅ | -
-| [**Reinforcement Learning for Control**](./reinforcement_learning)                    | - | - | coming soon.
+| [**ControlNet**](./controlnet) | ✅ | ✅ | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/controlnet.ipynb)
+| [**InstructPix2Pix**](./instruct_pix2pix) | ✅ | ✅ | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/InstructPix2Pix_using_diffusers.ipynb)
+| [**Reinforcement Learning for Control**](./reinforcement_learning)                    | - | - | [Notebook1](https://github.com/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_for_control.ipynb),  [Notebook2](https://github.com/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb)
 
 ## Community
 

examples/advanced_diffusion_training/train_dreambooth_lora_flux_advanced.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/advanced_diffusion_training/train_dreambooth_lora_sdxl_advanced.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/amused/train_amused.py

@@ -1,5 +1,5 @@
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team.
+# Copyright 2025 The HuggingFace Inc. team.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/community/README.md

@@ -24,32 +24,35 @@ Please also check out our [Community Scripts](https://github.com/huggingface/dif
 | Speech to Image                                                                                                                       | Using automatic-speech-recognition to transcribe text and Stable Diffusion to generate images                                                                                                                                                                                                                                                                                                                                                                                                            | [Speech to Image](#speech-to-image)                                                       |[Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/speech_to_image.ipynb)                                                                                                                                                                                                   |             [Mikail Duzenli](https://github.com/MikailINTech)
 | Wild Card Stable Diffusion                                                                                                            | Stable Diffusion Pipeline that supports prompts that contain wildcard terms (indicated by surrounding double underscores), with values instantiated randomly from a corresponding txt file or a dictionary of possible values                                                                                                                                                                                                                                                                            | [Wildcard Stable Diffusion](#wildcard-stable-diffusion)                                   | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/wildcard_stable_diffusion.ipynb)                                                                                                                                                                                 |              [Shyam Sudhakaran](https://github.com/shyamsn97) |
 | [Composable Stable Diffusion](https://energy-based-model.github.io/Compositional-Visual-Generation-with-Composable-Diffusion-Models/) | Stable Diffusion Pipeline that supports prompts that contain "|" in prompts (as an AND condition) and weights (separated by "|" as well) to positively / negatively weight prompts.                                                                                                                                                                                                                                                                                                            | [Composable Stable Diffusion](#composable-stable-diffusion)                               | -                                                                                                                                                                                                                  |                      [Mark Rich](https://github.com/MarkRich) |
-| Seed Resizing Stable Diffusion                                                                                                        | Stable Diffusion Pipeline that supports resizing an image and retaining the concepts of the 512 by 512 generation.                                                                                                                                                                                                                                                                                                                                                                                       | [Seed Resizing](#seed-resizing)                                                           | -                                                                                                                                                                                                                  |                      [Mark Rich](https://github.com/MarkRich) |
-| Imagic Stable Diffusion                                                                                                               | Stable Diffusion Pipeline that enables writing a text prompt to edit an existing image                                                                                                                                                                                                                                                                                                                                                                                                                   | [Imagic Stable Diffusion](#imagic-stable-diffusion)                                       | -                                                                                                                                                                                                                  |                      [Mark Rich](https://github.com/MarkRich) |
+| Seed Resizing Stable Diffusion                                                                                                        | Stable Diffusion Pipeline that supports resizing an image and retaining the concepts of the 512 by 512 generation.                                                                                                                                                                                                                                                                                                                                                                                       | [Seed Resizing](#seed-resizing)                                                           | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/seed_resizing.ipynb)                                                                                                                                                                                                                  |                      [Mark Rich](https://github.com/MarkRich) |
+| Imagic Stable Diffusion                                                                                                               | Stable Diffusion Pipeline that enables writing a text prompt to edit an existing image                                                                                                                                                                                                                                                                                                                                                                                                                   | [Imagic Stable Diffusion](#imagic-stable-diffusion)                                       | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/imagic_stable_diffusion.ipynb)                                                                                                                                                                                                  |                      [Mark Rich](https://github.com/MarkRich) |
 | Multilingual Stable Diffusion                                                                                                         | Stable Diffusion Pipeline that supports prompts in 50 different languages.                                                                                                                                                                                                                                                                                                                                                                                                                               | [Multilingual Stable Diffusion](#multilingual-stable-diffusion-pipeline)                  | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/multilingual_stable_diffusion.ipynb)                                                                                                                                                                             |          [Juan Carlos Piñeros](https://github.com/juancopi81) |
-| GlueGen Stable Diffusion                                                                                                         | Stable Diffusion Pipeline that supports prompts in different languages using GlueGen adapter.                                                                                                                                                                                                                                                                                                                                                                                                                               | [GlueGen Stable Diffusion](#gluegen-stable-diffusion-pipeline)                  | -                                                                                                                                                                                                                  |          [Phạm Hồng Vinh](https://github.com/rootonchair) |
+| GlueGen Stable Diffusion                                                                                                         | Stable Diffusion Pipeline that supports prompts in different languages using GlueGen adapter.                                                                                                                                                                                                                                                                                                                                                                                                                               | [GlueGen Stable Diffusion](#gluegen-stable-diffusion-pipeline)                  | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/gluegen_stable_diffusion.ipynb)                                                                                                                                                                                                                  |          [Phạm Hồng Vinh](https://github.com/rootonchair) |
 | Image to Image Inpainting Stable Diffusion                                                                                            | Stable Diffusion Pipeline that enables the overlaying of two images and subsequent inpainting                                                                                                                                                                                                                                                                                                                                                                                                            | [Image to Image Inpainting Stable Diffusion](#image-to-image-inpainting-stable-diffusion) | -                                                                                                                                                                                                                  |                    [Alex McKinney](https://github.com/vvvm23) |
-| Text Based Inpainting Stable Diffusion                                                                                                | Stable Diffusion Inpainting Pipeline that enables passing a text prompt to generate the mask for inpainting                                                                                                                                                                                                                                                                                                                                                                                              | [Text Based Inpainting Stable Diffusion](#text-based-inpainting-stable-diffusion)     | -                                                                                                                                                                                                                  |                   [Dhruv Karan](https://github.com/unography) |
+| Text Based Inpainting Stable Diffusion                                                                                                | Stable Diffusion Inpainting Pipeline that enables passing a text prompt to generate the mask for inpainting                                                                                                                                                                                                                                                                                                                                                                                              | [Text Based Inpainting Stable Diffusion](#text-based-inpainting-stable-diffusion)     | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/text_based_inpainting_stable_dffusion.ipynb)                                                                                                                                                                                                    |                   [Dhruv Karan](https://github.com/unography) |
 | Bit Diffusion                                                                                                                         | Diffusion on discrete data                                                                                                                                                                                                                                                                                                                                                                                                                                                                               | [Bit Diffusion](#bit-diffusion)                                                           | -  |                       [Stuti R.](https://github.com/kingstut) |
 | K-Diffusion Stable Diffusion                                                                                                          | Run Stable Diffusion with any of [K-Diffusion's samplers](https://github.com/crowsonkb/k-diffusion/blob/master/k_diffusion/sampling.py)                                                                                                                                                                                                                                                                                                                                                                  | [Stable Diffusion with K Diffusion](#stable-diffusion-with-k-diffusion)                   | -  |    [Patrick von Platen](https://github.com/patrickvonplaten/) |
 | Checkpoint Merger Pipeline                                                                                                            | Diffusion Pipeline that enables merging of saved model checkpoints                                                                                                                                                                                                                                                                                                                                                                                                                                       | [Checkpoint Merger Pipeline](#checkpoint-merger-pipeline)                                 | -                                                                                                                                                                                                                  | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) |
 | Stable Diffusion v1.1-1.4 Comparison                                                                                                  | Run all 4 model checkpoints for Stable Diffusion and compare their results together                                                                                                                                                                                                                                                                                                                                                                                                                      | [Stable Diffusion Comparison](#stable-diffusion-comparisons)                              | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/stable_diffusion_comparison.ipynb) |        [Suvaditya Mukherjee](https://github.com/suvadityamuk) |
-| MagicMix                                                                                                                              | Diffusion Pipeline for semantic mixing of an image and a text prompt                                                                                                                                                                                                                                                                                                                                                                                                                                     | [MagicMix](#magic-mix)                                                                    | - |                    [Partho Das](https://github.com/daspartho) |
+| MagicMix                                                                                                                              | Diffusion Pipeline for semantic mixing of an image and a text prompt                                                                                                                                                                                                                                                                                                                                                                                                                                     | [MagicMix](#magic-mix)                                                                    | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/magic_mix.ipynb) |                    [Partho Das](https://github.com/daspartho) |
 | Stable UnCLIP                                                                                                                         | Diffusion Pipeline for combining prior model (generate clip image embedding from text, UnCLIPPipeline `"kakaobrain/karlo-v1-alpha"`) and decoder pipeline (decode clip image embedding to image, StableDiffusionImageVariationPipeline `"lambdalabs/sd-image-variations-diffusers"` ).                                                                                                                                                                                                                   | [Stable UnCLIP](#stable-unclip)                                                           | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/stable_unclip.ipynb)  |                                [Ray Wang](https://wrong.wang) |
 | UnCLIP Text Interpolation Pipeline                                                                                                    | Diffusion Pipeline that allows passing two prompts and produces images while interpolating between the text-embeddings of the two prompts                                                                                                                                                                                                                                                                                                                                                                | [UnCLIP Text Interpolation Pipeline](#unclip-text-interpolation-pipeline)                 | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/unclip_text_interpolation.ipynb)| [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) |
-| UnCLIP Image Interpolation Pipeline                                                                                                   | Diffusion Pipeline that allows passing two images/image_embeddings and produces images while interpolating between their image-embeddings                                                                                                                                                                                                                                                                                                                                                                | [UnCLIP Image Interpolation Pipeline](#unclip-image-interpolation-pipeline)               | -                                                                                                                                                                                                                  | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) |
+| UnCLIP Image Interpolation Pipeline                                                                                                   | Diffusion Pipeline that allows passing two images/image_embeddings and produces images while interpolating between their image-embeddings                                                                                                                                                                                                                                                                                                                                                                | [UnCLIP Image Interpolation Pipeline](#unclip-image-interpolation-pipeline)               | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/unclip_image_interpolation.ipynb)| [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) |
 | DDIM Noise Comparative Analysis Pipeline                                                                                              | Investigating how the diffusion models learn visual concepts from each noise level (which is a contribution of [P2 weighting (CVPR 2022)](https://arxiv.org/abs/2204.00227))                                                                                                                                                                                                                                                                                                                             | [DDIM Noise Comparative Analysis Pipeline](#ddim-noise-comparative-analysis-pipeline)     | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/ddim_noise_comparative_analysis.ipynb)|              [Aengus (Duc-Anh)](https://github.com/aengusng8) |
-| CLIP Guided Img2Img Stable Diffusion Pipeline                                                                                         | Doing CLIP guidance for image to image generation with Stable Diffusion                                                                                                                                                                                                                                                                                                                                                                                                                                  | [CLIP Guided Img2Img Stable Diffusion](#clip-guided-img2img-stable-diffusion)             | - |               [Nipun Jindal](https://github.com/nipunjindal/) |
+| CLIP Guided Img2Img Stable Diffusion Pipeline                                                                                         | Doing CLIP guidance for image to image generation with Stable Diffusion                                                                                                                                                                                                                                                                                                                                                                                                                                  | [CLIP Guided Img2Img Stable Diffusion](#clip-guided-img2img-stable-diffusion)             | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/clip_guided_img2img_stable_diffusion.ipynb) |               [Nipun Jindal](https://github.com/nipunjindal/) |
 | TensorRT Stable Diffusion Text to Image Pipeline                                                                                                    | Accelerates the Stable Diffusion Text2Image Pipeline using TensorRT                                                                                                                                                                                                                                                                                                                                                                                                                                      | [TensorRT Stable Diffusion Text to Image Pipeline](#tensorrt-text2image-stable-diffusion-pipeline)      | - |              [Asfiya Baig](https://github.com/asfiyab-nvidia) |
 | EDICT Image Editing Pipeline                                                                                                          | Diffusion pipeline for text-guided image editing                                                                                                                                                                                                                                                                                                                                                                                                                                                         | [EDICT Image Editing Pipeline](#edict-image-editing-pipeline)                             | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/edict_image_pipeline.ipynb) |                    [Joqsan Azocar](https://github.com/Joqsan) |
 | Stable Diffusion RePaint                                                                                                              | Stable Diffusion pipeline using [RePaint](https://arxiv.org/abs/2201.09865) for inpainting.                                                                                                                                                                                                                                                                                                                                                                                                               | [Stable Diffusion RePaint](#stable-diffusion-repaint )|[Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/stable_diffusion_repaint.ipynb)|                  [Markus Pobitzer](https://github.com/Markus-Pobitzer) |
 | TensorRT Stable Diffusion Image to Image Pipeline                                                                                                    | Accelerates the Stable Diffusion Image2Image Pipeline using TensorRT                                                                                                                                                                                                                                                                                                                                                                                                                                      | [TensorRT Stable Diffusion Image to Image Pipeline](#tensorrt-image2image-stable-diffusion-pipeline)      | - |              [Asfiya Baig](https://github.com/asfiyab-nvidia) |
 | Stable Diffusion IPEX Pipeline | Accelerate Stable Diffusion inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [Stable Diffusion on IPEX](#stable-diffusion-on-ipex) | - | [Yingjie Han](https://github.com/yingjie-han/) |
-| CLIP Guided Images Mixing Stable Diffusion Pipeline | Сombine images using usual diffusion models. | [CLIP Guided Images Mixing Using Stable Diffusion](#clip-guided-images-mixing-with-stable-diffusion) | - | [Karachev Denis](https://github.com/TheDenk) |
+| CLIP Guided Images Mixing Stable Diffusion Pipeline | Сombine images using usual diffusion models. | [CLIP Guided Images Mixing Using Stable Diffusion](#clip-guided-images-mixing-with-stable-diffusion) | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/clip_guided_images_mixing_with_stable_diffusion.ipynb) | [Karachev Denis](https://github.com/TheDenk) |
 | TensorRT Stable Diffusion Inpainting Pipeline                                                                                                    | Accelerates the Stable Diffusion Inpainting Pipeline using TensorRT                                                                                                                                                                                                                                                                                                                                                                                                                                      | [TensorRT Stable Diffusion Inpainting Pipeline](#tensorrt-inpainting-stable-diffusion-pipeline)      | - |              [Asfiya Baig](https://github.com/asfiyab-nvidia) |
 |   IADB Pipeline                                                                                                    | Implementation of [Iterative α-(de)Blending: a Minimalist Deterministic Diffusion Model](https://arxiv.org/abs/2305.03486)                                                                                                                                                                                                                                                                                                                                                                                                                                      | [IADB Pipeline](#iadb-pipeline)      | - |              [Thomas Chambon](https://github.com/tchambon)
 |   Zero1to3 Pipeline                                                                                                    | Implementation of [Zero-1-to-3: Zero-shot One Image to 3D Object](https://arxiv.org/abs/2303.11328)                                                                                                                                                                                                                                                                                                                                                                                                                                      | [Zero1to3 Pipeline](#zero1to3-pipeline)      | - |              [Xin Kong](https://github.com/kxhit) |
 | Stable Diffusion XL Long Weighted Prompt Pipeline | A pipeline support unlimited length of prompt and negative prompt, use A1111 style of prompt weighting | [Stable Diffusion XL Long Weighted Prompt Pipeline](#stable-diffusion-xl-long-weighted-prompt-pipeline) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1LsqilswLR40XLLcp6XFOl5nKb_wOe26W?usp=sharing) | [Andrew Zhu](https://xhinker.medium.com/) |
+| Stable Diffusion Mixture Tiling Pipeline SD 1.5 | A pipeline generates cohesive images by integrating multiple diffusion processes, each focused on a specific image region and considering boundary effects for smooth blending | [Stable Diffusion Mixture Tiling Pipeline SD 1.5](#stable-diffusion-mixture-tiling-pipeline-sd-15) | [![Hugging Face Space](https://img.shields.io/badge/🤗%20Hugging%20Face-Space-yellow)](https://huggingface.co/spaces/albarji/mixture-of-diffusers) | [Álvaro B Jiménez](https://github.com/albarji/) |
+| Stable Diffusion Mixture Canvas Pipeline SD 1.5 | A pipeline generates cohesive images by integrating multiple diffusion processes, each focused on a specific image region and considering boundary effects for smooth blending. Works by defining a list of Text2Image region objects that detail the region of influence of each diffuser. | [Stable Diffusion Mixture Canvas Pipeline SD 1.5](#stable-diffusion-mixture-canvas-pipeline-sd-15) | [![Hugging Face Space](https://img.shields.io/badge/🤗%20Hugging%20Face-Space-yellow)](https://huggingface.co/spaces/albarji/mixture-of-diffusers) | [Álvaro B Jiménez](https://github.com/albarji/) |
+| Stable Diffusion Mixture Tiling Pipeline SDXL | A pipeline generates cohesive images by integrating multiple diffusion processes, each focused on a specific image region and considering boundary effects for smooth blending | [Stable Diffusion Mixture Tiling Pipeline SDXL](#stable-diffusion-mixture-tiling-pipeline-sdxl) | [![Hugging Face Space](https://img.shields.io/badge/🤗%20Hugging%20Face-Space-yellow)](https://huggingface.co/spaces/elismasilva/mixture-of-diffusers-sdxl-tiling) | [Eliseu Silva](https://github.com/DEVAIEXP/) |
 | FABRIC - Stable Diffusion with feedback Pipeline | pipeline supports feedback from liked and disliked images | [Stable Diffusion Fabric Pipeline](#stable-diffusion-fabric-pipeline) | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/stable_diffusion_fabric.ipynb)| [Shauray Singh](https://shauray8.github.io/about_shauray/) |
 | sketch inpaint - Inpainting with non-inpaint Stable Diffusion | sketch inpaint much like in automatic1111 | [Masked Im2Im Stable Diffusion Pipeline](#stable-diffusion-masked-im2im) | - | [Anatoly Belikov](https://github.com/noskill) |
 | sketch inpaint xl - Inpainting with non-inpaint Stable Diffusion | sketch inpaint much like in automatic1111 | [Masked Im2Im Stable Diffusion XL Pipeline](#stable-diffusion-xl-masked-im2im) | - | [Anatoly Belikov](https://github.com/noskill) |
@@ -57,7 +60,7 @@ Please also check out our [Community Scripts](https://github.com/huggingface/dif
 |   Latent Consistency Pipeline                                                                                                    | Implementation of [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://arxiv.org/abs/2310.04378)                                                                                                                                                                                                                                                                                                                                                                                                                                      | [Latent Consistency Pipeline](#latent-consistency-pipeline)      | - |              [Simian Luo](https://github.com/luosiallen) |
 |   Latent Consistency Img2img Pipeline                                                                                                    | Img2img pipeline for Latent Consistency Models                                                                                                                                                                                                                                                                                                                                                                                                                                    | [Latent Consistency Img2Img Pipeline](#latent-consistency-img2img-pipeline)      | - |              [Logan Zoellner](https://github.com/nagolinc) |
 |   Latent Consistency Interpolation Pipeline                                                                                                    | Interpolate the latent space of Latent Consistency Models with multiple prompts                                                                                                                                                                                                                                                                                                                                                                                                                                    | [Latent Consistency Interpolation Pipeline](#latent-consistency-interpolation-pipeline) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1pK3NrLWJSiJsBynLns1K1-IDTW9zbPvl?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) |
-| SDE Drag Pipeline                                                                                                                         | The pipeline supports drag editing of images using stochastic differential equations                                                                                                                                                                                                                                                                                                                                                                                                                | [SDE Drag Pipeline](#sde-drag-pipeline)                                                     | - | [NieShen](https://github.com/NieShenRuc) [Fengqi Zhu](https://github.com/Monohydroxides) |
+| SDE Drag Pipeline                                                                                                                         | The pipeline supports drag editing of images using stochastic differential equations                                                                                                                                                                                                                                                                                                                                                                                                                | [SDE Drag Pipeline](#sde-drag-pipeline)                                                     | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/sde_drag.ipynb) | [NieShen](https://github.com/NieShenRuc) [Fengqi Zhu](https://github.com/Monohydroxides) |
 |   Regional Prompting Pipeline                                                                                               | Assign multiple prompts for different regions                                                                                                                                                                                                                                                                                                                                                    |  [Regional Prompting Pipeline](#regional-prompting-pipeline) | - | [hako-mikan](https://github.com/hako-mikan) |
 | LDM3D-sr (LDM3D upscaler)                                                                                                             | Upscale low resolution RGB and depth inputs to high resolution                                                                                                                                                                                                                                                                                                                                                                                                                              | [StableDiffusionUpscaleLDM3D Pipeline](https://github.com/estelleafl/diffusers/tree/ldm3d_upscaler_community/examples/community#stablediffusionupscaleldm3d-pipeline)                                                                             | -                                                                                                                                                                                                             |                                                        [Estelle Aflalo](https://github.com/estelleafl) |
 | AnimateDiff ControlNet Pipeline                                                                                                    | Combines AnimateDiff with precise motion control using ControlNets                                                                                                                                                                                                                                                                                                                                                                                                                                    | [AnimateDiff ControlNet Pipeline](#animatediff-controlnet-pipeline) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1SKboYeGjEQmQPWoFC0aLYpBlYdHXkvAu?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) and [Edoardo Botta](https://github.com/EdoardoBotta) |
@@ -77,6 +80,8 @@ Please also check out our [Community Scripts](https://github.com/huggingface/dif
 PIXART-α Controlnet pipeline | Implementation of the controlnet model for pixart alpha and its diffusers pipeline | [PIXART-α Controlnet pipeline](#pixart-α-controlnet-pipeline) | - | [Raul Ciotescu](https://github.com/raulc0399/) |
 | HunyuanDiT Differential Diffusion Pipeline | Applies [Differential Diffusion](https://github.com/exx8/differential-diffusion) to [HunyuanDiT](https://github.com/huggingface/diffusers/pull/8240). | [HunyuanDiT with Differential Diffusion](#hunyuandit-with-differential-diffusion) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing) | [Monjoy Choudhury](https://github.com/MnCSSJ4x) |
 | [🪆Matryoshka Diffusion Models](https://huggingface.co/papers/2310.15111) | A diffusion process that denoises inputs at multiple resolutions jointly and uses a NestedUNet architecture where features and parameters for small scale inputs are nested within those of the large scales. See [original codebase](https://github.com/apple/ml-mdm). | [🪆Matryoshka Diffusion Models](#matryoshka-diffusion-models) | [![Hugging Face Space](https://img.shields.io/badge/🤗%20Hugging%20Face-Space-yellow)](https://huggingface.co/spaces/pcuenq/mdm) [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/gist/tolgacangoz/1f54875fc7aeaabcf284ebde64820966/matryoshka_hf.ipynb) | [M. Tolga Cangöz](https://github.com/tolgacangoz) |
+| Stable Diffusion XL Attentive Eraser Pipeline |[[AAAI2025 Oral] Attentive Eraser](https://github.com/Anonym0u3/AttentiveEraser) is a novel tuning-free method that enhances object removal capabilities in pre-trained diffusion models.|[Stable Diffusion XL Attentive Eraser Pipeline](#stable-diffusion-xl-attentive-eraser-pipeline)|-|[Wenhao Sun](https://github.com/Anonym0u3) and [Benlei Cui](https://github.com/Benny079)|
+| Perturbed-Attention Guidance |StableDiffusionPAGPipeline is a modification of StableDiffusionPipeline to support Perturbed-Attention Guidance (PAG).|[Perturbed-Attention Guidance](#perturbed-attention-guidance)|[Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/perturbed_attention_guidance.ipynb)|[Hyoungwon Cho](https://github.com/HyoungwonCho)|
 
 To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly.
 
@@ -947,10 +952,15 @@ image.save('./imagic/imagic_image_alpha_2.png')
 Test seed resizing. Originally generate an image in 512 by 512, then generate image with same seed at 512 by 592 using seed resizing. Finally, generate 512 by 592 using original stable diffusion pipeline.
 
 ```python
+import os
 import torch as th
 import numpy as np
 from diffusers import DiffusionPipeline
 
+# Ensure the save directory exists or create it
+save_dir = './seed_resize/'
+os.makedirs(save_dir, exist_ok=True)
+
 has_cuda = th.cuda.is_available()
 device = th.device('cpu' if not has_cuda else 'cuda')
 
@@ -964,7 +974,6 @@ def dummy(images, **kwargs):
 
 pipe.safety_checker = dummy
 
-
 images = []
 th.manual_seed(0)
 generator = th.Generator("cuda").manual_seed(0)
@@ -983,15 +992,14 @@ res = pipe(
     width=width,
     generator=generator)
 image = res.images[0]
-image.save('./seed_resize/seed_resize_{w}_{h}_image.png'.format(w=width, h=height))
-
+image.save(os.path.join(save_dir, 'seed_resize_{w}_{h}_image.png'.format(w=width, h=height)))
 
 th.manual_seed(0)
 generator = th.Generator("cuda").manual_seed(0)
 
 pipe = DiffusionPipeline.from_pretrained(
     "CompVis/stable-diffusion-v1-4",
-    custom_pipeline="/home/mark/open_source/diffusers/examples/community/"
+    custom_pipeline="seed_resize_stable_diffusion"
 ).to(device)
 
 width = 512
@@ -1005,11 +1013,11 @@ res = pipe(
     width=width,
     generator=generator)
 image = res.images[0]
-image.save('./seed_resize/seed_resize_{w}_{h}_image.png'.format(w=width, h=height))
+image.save(os.path.join(save_dir, 'seed_resize_{w}_{h}_image.png'.format(w=width, h=height)))
 
 pipe_compare = DiffusionPipeline.from_pretrained(
     "CompVis/stable-diffusion-v1-4",
-    custom_pipeline="/home/mark/open_source/diffusers/examples/community/"
+    custom_pipeline="seed_resize_stable_diffusion"
 ).to(device)
 
 res = pipe_compare(
@@ -1022,7 +1030,7 @@ res = pipe_compare(
 )
 
 image = res.images[0]
-image.save('./seed_resize/seed_resize_{w}_{h}_image_compare.png'.format(w=width, h=height))
+image.save(os.path.join(save_dir, 'seed_resize_{w}_{h}_image_compare.png'.format(w=width, h=height)))
 ```
 
 ### Multilingual Stable Diffusion Pipeline
@@ -1099,38 +1107,100 @@ GlueGen is a minimal adapter that allows alignment between any encoder (Text Enc
 Make sure you downloaded `gluenet_French_clip_overnorm_over3_noln.ckpt` for French (there are also pre-trained weights for Chinese, Italian, Japanese, Spanish or train your own) at [GlueGen's official repo](https://github.com/salesforce/GlueGen/tree/main).
 
 ```python
-from PIL import Image
-
+import os
+import gc
+import urllib.request
 import torch
-
-from transformers import AutoModel, AutoTokenizer
-
+from transformers import XLMRobertaTokenizer, XLMRobertaForMaskedLM, CLIPTokenizer, CLIPTextModel
 from diffusers import DiffusionPipeline
 
-if __name__ == "__main__":
-    device = "cuda"
+# Download checkpoints
+CHECKPOINTS = [
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_Chinese_clip_overnorm_over3_noln.ckpt",
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_French_clip_overnorm_over3_noln.ckpt",
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_Italian_clip_overnorm_over3_noln.ckpt",
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_Japanese_clip_overnorm_over3_noln.ckpt",
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_Spanish_clip_overnorm_over3_noln.ckpt",
+    "https://storage.googleapis.com/sfr-gluegen-data-research/checkpoints_all/gluenet_checkpoint/gluenet_sound2img_audioclip_us8k.ckpt"
+]
 
-    lm_model_id = "xlm-roberta-large"
-    token_max_length = 77
+LANGUAGE_PROMPTS = {
+    "French": "une voiture sur la plage",
+    #"Chinese": "海滩上的一辆车",
+    #"Italian": "una macchina sulla spiaggia",
+    #"Japanese": "浜辺の車",
+    #"Spanish": "un coche en la playa"
+}
 
-    text_encoder = AutoModel.from_pretrained(lm_model_id)
-    tokenizer = AutoTokenizer.from_pretrained(lm_model_id, model_max_length=token_max_length, use_fast=False)
+def download_checkpoints(checkpoint_dir):
+    os.makedirs(checkpoint_dir, exist_ok=True)
+    for url in CHECKPOINTS:
+        filename = os.path.join(checkpoint_dir, os.path.basename(url))
+        if not os.path.exists(filename):
+            print(f"Downloading {filename}...")
+            urllib.request.urlretrieve(url, filename)
+            print(f"Downloaded {filename}")
+        else:
+            print(f"Checkpoint {filename} already exists, skipping download.")
+    return checkpoint_dir
 
-    tensor_norm = torch.Tensor([[43.8203],[28.3668],[27.9345],[28.0084],[28.2958],[28.2576],[28.3373],[28.2695],[28.4097],[28.2790],[28.2825],[28.2807],[28.2775],[28.2708],[28.2682],[28.2624],[28.2589],[28.2611],[28.2616],[28.2639],[28.2613],[28.2566],[28.2615],[28.2665],[28.2799],[28.2885],[28.2852],[28.2863],[28.2780],[28.2818],[28.2764],[28.2532],[28.2412],[28.2336],[28.2514],[28.2734],[28.2763],[28.2977],[28.2971],[28.2948],[28.2818],[28.2676],[28.2831],[28.2890],[28.2979],[28.2999],[28.3117],[28.3363],[28.3554],[28.3626],[28.3589],[28.3597],[28.3543],[28.3660],[28.3731],[28.3717],[28.3812],[28.3753],[28.3810],[28.3777],[28.3693],[28.3713],[28.3670],[28.3691],[28.3679],[28.3624],[28.3703],[28.3703],[28.3720],[28.3594],[28.3576],[28.3562],[28.3438],[28.3376],[28.3389],[28.3433],[28.3191]])
+def load_checkpoint(pipeline, checkpoint_path, device):
+    state_dict = torch.load(checkpoint_path, map_location=device)
+    state_dict = state_dict.get("state_dict", state_dict)
+    missing_keys, unexpected_keys = pipeline.unet.load_state_dict(state_dict, strict=False)
+    return pipeline
 
-    pipeline = DiffusionPipeline.from_pretrained(
-        "stable-diffusion-v1-5/stable-diffusion-v1-5",
-        text_encoder=text_encoder,
-        tokenizer=tokenizer,
-        custom_pipeline="gluegen"
-    ).to(device)
-    pipeline.load_language_adapter("gluenet_French_clip_overnorm_over3_noln.ckpt", num_token=token_max_length, dim=1024, dim_out=768, tensor_norm=tensor_norm)
+def generate_image(pipeline, prompt, device, output_path):
+    with torch.inference_mode():
+        image = pipeline(
+            prompt,
+            generator=torch.Generator(device=device).manual_seed(42),
+            num_inference_steps=50
+        ).images[0]
+        image.save(output_path)
+        print(f"Image saved to {output_path}")
 
-    prompt = "une voiture sur la plage"
+checkpoint_dir = download_checkpoints("./checkpoints_all/gluenet_checkpoint")
+device = "cuda" if torch.cuda.is_available() else "cpu"
+print(f"Using device: {device}")
 
-    generator = torch.Generator(device=device).manual_seed(42)
-    image = pipeline(prompt, generator=generator).images[0]
-    image.save("gluegen_output_fr.png")
+tokenizer = XLMRobertaTokenizer.from_pretrained("xlm-roberta-base", use_fast=False)
+model = XLMRobertaForMaskedLM.from_pretrained("xlm-roberta-base").to(device)
+inputs = tokenizer("Ceci est une phrase incomplète avec un [MASK].", return_tensors="pt").to(device)
+with torch.inference_mode():
+    _ = model(**inputs)
+
+
+clip_tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14")
+clip_text_encoder = CLIPTextModel.from_pretrained("openai/clip-vit-large-patch14").to(device)
+
+# Initialize pipeline
+pipeline = DiffusionPipeline.from_pretrained(
+    "stable-diffusion-v1-5/stable-diffusion-v1-5",
+    text_encoder=clip_text_encoder,
+    tokenizer=clip_tokenizer,
+    custom_pipeline="gluegen",
+    safety_checker=None
+).to(device)
+
+os.makedirs("outputs", exist_ok=True)
+
+# Generate images
+for language, prompt in LANGUAGE_PROMPTS.items():
+
+    checkpoint_file = f"gluenet_{language}_clip_overnorm_over3_noln.ckpt"
+    checkpoint_path = os.path.join(checkpoint_dir, checkpoint_file)
+    try:
+        pipeline = load_checkpoint(pipeline, checkpoint_path, device)
+        output_path = f"outputs/gluegen_output_{language.lower()}.png"
+        generate_image(pipeline, prompt, device, output_path)
+    except Exception as e:
+        print(f"Error processing {language} model: {e}")
+        continue
+
+    if torch.cuda.is_available():
+        torch.cuda.empty_cache()
+    gc.collect()
 ```
 
 Which will produce:
@@ -1181,28 +1251,49 @@ Currently uses the CLIPSeg model for mask generation, then calls the standard St
 ```python
 from transformers import CLIPSegProcessor, CLIPSegForImageSegmentation
 from diffusers import DiffusionPipeline
-
 from PIL import Image
 import requests
+import torch
 
+# Load CLIPSeg model and processor
 processor = CLIPSegProcessor.from_pretrained("CIDAS/clipseg-rd64-refined")
-model = CLIPSegForImageSegmentation.from_pretrained("CIDAS/clipseg-rd64-refined")
+model = CLIPSegForImageSegmentation.from_pretrained("CIDAS/clipseg-rd64-refined").to("cuda")
 
+# Load Stable Diffusion Inpainting Pipeline with custom pipeline
 pipe = DiffusionPipeline.from_pretrained(
     "runwayml/stable-diffusion-inpainting",
     custom_pipeline="text_inpainting",
     segmentation_model=model,
     segmentation_processor=processor
-)
-pipe = pipe.to("cuda")
-
+).to("cuda")
 
+# Load input image
 url = "https://github.com/timojl/clipseg/blob/master/example_image.jpg?raw=true"
-image = Image.open(requests.get(url, stream=True).raw).resize((512, 512))
-text = "a glass"  # will mask out this text
-prompt = "a cup"  # the masked out region will be replaced with this
+image = Image.open(requests.get(url, stream=True).raw)
 
-image = pipe(image=image, text=text, prompt=prompt).images[0]
+# Step 1: Resize input image for CLIPSeg (224x224)
+segmentation_input = image.resize((224, 224))
+
+# Step 2: Generate segmentation mask
+text = "a glass"  # Object to mask
+inputs = processor(text=text, images=segmentation_input, return_tensors="pt").to("cuda")
+
+with torch.no_grad():
+    mask = model(**inputs).logits.sigmoid()  # Get segmentation mask
+
+# Resize mask back to 512x512 for SD inpainting
+mask = torch.nn.functional.interpolate(mask.unsqueeze(0), size=(512, 512), mode="bilinear").squeeze(0)
+
+# Step 3: Resize input image for Stable Diffusion
+image = image.resize((512, 512))
+
+# Step 4: Run inpainting with Stable Diffusion
+prompt = "a cup"  # The masked-out region will be replaced with this
+result = pipe(image=image, mask=mask, prompt=prompt,text=text).images[0]
+
+# Save output
+result.save("inpainting_output.png")
+print("Inpainting completed. Image saved as 'inpainting_output.png'.")
 ```
 
 ### Bit Diffusion
@@ -1378,8 +1469,10 @@ There are 3 parameters for the method-
 Here is an example usage-
 
 ```python
+import requests
 from diffusers import DiffusionPipeline, DDIMScheduler
 from PIL import Image
+from io import BytesIO
 
 pipe = DiffusionPipeline.from_pretrained(
     "CompVis/stable-diffusion-v1-4",
@@ -1387,9 +1480,11 @@ pipe = DiffusionPipeline.from_pretrained(
     scheduler=DDIMScheduler.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="scheduler"),
 ).to('cuda')
 
-img = Image.open('phone.jpg')
+url = "https://user-images.githubusercontent.com/59410571/209578593-141467c7-d831-4792-8b9a-b17dc5e47816.jpg"
+response = requests.get(url)
+image = Image.open(BytesIO(response.content)).convert("RGB")  # Convert to RGB to avoid issues
 mix_img = pipe(
-    img,
+    image,
     prompt='bed',
     kmin=0.3,
     kmax=0.5,
@@ -1542,6 +1637,8 @@ This Diffusion Pipeline takes two images or an image_embeddings tensor of size 2
 import torch
 from diffusers import DiffusionPipeline
 from PIL import Image
+import requests
+from io import BytesIO
 
 device = torch.device("cpu" if not torch.cuda.is_available() else "cuda")
 dtype = torch.float16 if torch.cuda.is_available() else torch.bfloat16
@@ -1553,13 +1650,25 @@ pipe = DiffusionPipeline.from_pretrained(
 )
 pipe.to(device)
 
-images = [Image.open('./starry_night.jpg'), Image.open('./flowers.jpg')]
+# List of image URLs
+image_urls = [
+    'https://camo.githubusercontent.com/ef13c8059b12947c0d5e8d3ea88900de6bf1cd76bbf61ace3928e824c491290e/68747470733a2f2f68756767696e67666163652e636f2f64617461736574732f4e616761536169416268696e61792f556e434c4950496d616765496e746572706f6c6174696f6e53616d706c65732f7265736f6c76652f6d61696e2f7374617272795f6e696768742e6a7067',
+    'https://camo.githubusercontent.com/d1947ab7c49ae3f550c28409d5e8b120df48e456559cf4557306c0848337702c/68747470733a2f2f68756767696e67666163652e636f2f64617461736574732f4e616761536169416268696e61792f556e434c4950496d616765496e746572706f6c6174696f6e53616d706c65732f7265736f6c76652f6d61696e2f666c6f776572732e6a7067'
+]
+
+# Open images from URLs
+images = []
+for url in image_urls:
+    response = requests.get(url)
+    img = Image.open(BytesIO(response.content))
+    images.append(img)
+
 # For best results keep the prompts close in length to each other. Of course, feel free to try out with differing lengths.
 generator = torch.Generator(device=device).manual_seed(42)
 
 output = pipe(image=images, steps=6, generator=generator)
 
-for i,image in enumerate(output.images):
+for i, image in enumerate(output.images):
     image.save('starry_to_flowers_%s.jpg' % i)
 ```
 
@@ -1636,37 +1745,51 @@ from diffusers import DiffusionPipeline
 from PIL import Image
 from transformers import CLIPImageProcessor, CLIPModel
 
+# Load CLIP model and feature extractor
 feature_extractor = CLIPImageProcessor.from_pretrained(
     "laion/CLIP-ViT-B-32-laion2B-s34B-b79K"
 )
 clip_model = CLIPModel.from_pretrained(
     "laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16
 )
+
+# Load guided pipeline
 guided_pipeline = DiffusionPipeline.from_pretrained(
     "CompVis/stable-diffusion-v1-4",
-    # custom_pipeline="clip_guided_stable_diffusion",
-    custom_pipeline="/home/njindal/diffusers/examples/community/clip_guided_stable_diffusion.py",
+    custom_pipeline="clip_guided_stable_diffusion_img2img",
     clip_model=clip_model,
     feature_extractor=feature_extractor,
     torch_dtype=torch.float16,
 )
 guided_pipeline.enable_attention_slicing()
 guided_pipeline = guided_pipeline.to("cuda")
+
+# Define prompt and fetch image
 prompt = "fantasy book cover, full moon, fantasy forest landscape, golden vector elements, fantasy magic, dark light night, intricate, elegant, sharp focus, illustration, highly detailed, digital painting, concept art, matte, art by WLOP and Artgerm and Albert Bierstadt, masterpiece"
 url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
 response = requests.get(url)
-init_image = Image.open(BytesIO(response.content)).convert("RGB")
+edit_image = Image.open(BytesIO(response.content)).convert("RGB")
+
+# Run the pipeline
 image = guided_pipeline(
     prompt=prompt,
-    num_inference_steps=30,
-    image=init_image,
-    strength=0.75,
-    guidance_scale=7.5,
-    clip_guidance_scale=100,
-    num_cutouts=4,
-    use_cutouts=False,
+    height=512,  # Height of the output image
+    width=512,   # Width of the output image
+    image=edit_image,  # Input image to guide the diffusion
+    strength=0.75,  # How much to transform the input image
+    num_inference_steps=30,  # Number of diffusion steps
+    guidance_scale=7.5,  # Scale of the classifier-free guidance
+    clip_guidance_scale=100,  # Scale of the CLIP guidance
+    num_images_per_prompt=1,  # Generate one image per prompt
+    eta=0.0,  # Noise scheduling parameter
+    num_cutouts=4,  # Number of cutouts for CLIP guidance
+    use_cutouts=False,  # Whether to use cutouts
+    output_type="pil",  # Output as PIL image
 ).images[0]
-display(image)
+
+# Display the generated image
+image.show()
+
 ```
 
 Init Image
@@ -2243,6 +2366,85 @@ CLIP guided stable diffusion images mixing pipeline allows to combine two images
 This approach is using (optional) CoCa model to avoid writing image description.
 [More code examples](https://github.com/TheDenk/images_mixing)
 
+### Example Images Mixing (with CoCa)
+
+```python
+import PIL
+import torch
+import requests
+import open_clip
+from open_clip import SimpleTokenizer
+from io import BytesIO
+from diffusers import DiffusionPipeline
+from transformers import CLIPImageProcessor, CLIPModel
+
+
+def download_image(url):
+    response = requests.get(url)
+    return PIL.Image.open(BytesIO(response.content)).convert("RGB")
+
+# Loading additional models
+feature_extractor = CLIPImageProcessor.from_pretrained(
+    "laion/CLIP-ViT-B-32-laion2B-s34B-b79K"
+)
+clip_model = CLIPModel.from_pretrained(
+    "laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16
+)
+coca_model = open_clip.create_model('coca_ViT-L-14', pretrained='laion2B-s13B-b90k').to('cuda')
+coca_model.dtype = torch.float16
+coca_transform = open_clip.image_transform(
+    coca_model.visual.image_size,
+    is_train=False,
+    mean=getattr(coca_model.visual, 'image_mean', None),
+    std=getattr(coca_model.visual, 'image_std', None),
+)
+coca_tokenizer = SimpleTokenizer()
+
+# Pipeline creating
+mixing_pipeline = DiffusionPipeline.from_pretrained(
+    "CompVis/stable-diffusion-v1-4",
+    custom_pipeline="clip_guided_images_mixing_stable_diffusion",
+    clip_model=clip_model,
+    feature_extractor=feature_extractor,
+    coca_model=coca_model,
+    coca_tokenizer=coca_tokenizer,
+    coca_transform=coca_transform,
+    torch_dtype=torch.float16,
+)
+mixing_pipeline.enable_attention_slicing()
+mixing_pipeline = mixing_pipeline.to("cuda")
+
+# Pipeline running
+generator = torch.Generator(device="cuda").manual_seed(17)
+
+def download_image(url):
+    response = requests.get(url)
+    return PIL.Image.open(BytesIO(response.content)).convert("RGB")
+
+content_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir.jpg")
+style_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/gigachad.jpg")
+
+pipe_images = mixing_pipeline(
+    num_inference_steps=50,
+    content_image=content_image,
+    style_image=style_image,
+    noise_strength=0.65,
+    slerp_latent_style_strength=0.9,
+    slerp_prompt_style_strength=0.1,
+    slerp_clip_image_style_strength=0.1,
+    guidance_scale=9.0,
+    batch_size=1,
+    clip_guidance_scale=100,
+    generator=generator,
+).images
+
+output_path = "mixed_output.jpg"
+pipe_images[0].save(output_path)
+print(f"Image saved successfully at {output_path}")
+```
+
+![image_mixing_result](https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir_gigachad.png)
+
 ### Stable Diffusion XL Long Weighted Prompt Pipeline
 
 This SDXL pipeline supports unlimited length prompt and negative prompt, compatible with A1111 prompt weighted style.
@@ -2308,83 +2510,7 @@ In the above code, the `prompt2` is appended to the `prompt`, which is more than
 
 For more results, checkout [PR #6114](https://github.com/huggingface/diffusers/pull/6114).
 
-### Example Images Mixing (with CoCa)
-
-```python
-import requests
-from io import BytesIO
-
-import PIL
-import torch
-import open_clip
-from open_clip import SimpleTokenizer
-from diffusers import DiffusionPipeline
-from transformers import CLIPImageProcessor, CLIPModel
-
-
-def download_image(url):
-    response = requests.get(url)
-    return PIL.Image.open(BytesIO(response.content)).convert("RGB")
-
-# Loading additional models
-feature_extractor = CLIPImageProcessor.from_pretrained(
-    "laion/CLIP-ViT-B-32-laion2B-s34B-b79K"
-)
-clip_model = CLIPModel.from_pretrained(
-    "laion/CLIP-ViT-B-32-laion2B-s34B-b79K", torch_dtype=torch.float16
-)
-coca_model = open_clip.create_model('coca_ViT-L-14', pretrained='laion2B-s13B-b90k').to('cuda')
-coca_model.dtype = torch.float16
-coca_transform = open_clip.image_transform(
-    coca_model.visual.image_size,
-    is_train=False,
-    mean=getattr(coca_model.visual, 'image_mean', None),
-    std=getattr(coca_model.visual, 'image_std', None),
-)
-coca_tokenizer = SimpleTokenizer()
-
-# Pipeline creating
-mixing_pipeline = DiffusionPipeline.from_pretrained(
-    "CompVis/stable-diffusion-v1-4",
-    custom_pipeline="clip_guided_images_mixing_stable_diffusion",
-    clip_model=clip_model,
-    feature_extractor=feature_extractor,
-    coca_model=coca_model,
-    coca_tokenizer=coca_tokenizer,
-    coca_transform=coca_transform,
-    torch_dtype=torch.float16,
-)
-mixing_pipeline.enable_attention_slicing()
-mixing_pipeline = mixing_pipeline.to("cuda")
-
-# Pipeline running
-generator = torch.Generator(device="cuda").manual_seed(17)
-
-def download_image(url):
-    response = requests.get(url)
-    return PIL.Image.open(BytesIO(response.content)).convert("RGB")
-
-content_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir.jpg")
-style_image = download_image("https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/gigachad.jpg")
-
-pipe_images = mixing_pipeline(
-    num_inference_steps=50,
-    content_image=content_image,
-    style_image=style_image,
-    noise_strength=0.65,
-    slerp_latent_style_strength=0.9,
-    slerp_prompt_style_strength=0.1,
-    slerp_clip_image_style_strength=0.1,
-    guidance_scale=9.0,
-    batch_size=1,
-    clip_guidance_scale=100,
-    generator=generator,
-).images
-```
-
-![image_mixing_result](https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir_gigachad.png)
-
-### Stable Diffusion Mixture Tiling
+### Stable Diffusion Mixture Tiling Pipeline SD 1.5
 
 This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
 
@@ -2415,6 +2541,95 @@ image = pipeline(
 
 ![mixture_tiling_results](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/mixture_tiling.png)
 
+### Stable Diffusion Mixture Canvas Pipeline SD 1.5
+
+This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
+
+```python
+from PIL import Image
+from diffusers import LMSDiscreteScheduler, DiffusionPipeline
+from diffusers.pipelines.pipeline_utils import Image2ImageRegion, Text2ImageRegion, preprocess_image
+
+
+# Load and preprocess guide image
+iic_image = preprocess_image(Image.open("input_image.png").convert("RGB"))
+
+# Create scheduler and model (similar to StableDiffusionPipeline)
+scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
+pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler).to("cuda:0", custom_pipeline="mixture_canvas")
+pipeline.to("cuda")
+
+# Mixture of Diffusers generation
+output = pipeline(
+    canvas_height=800,
+    canvas_width=352,
+    regions=[
+        Text2ImageRegion(0, 800, 0, 352, guidance_scale=8,
+            prompt=f"best quality, masterpiece, WLOP, sakimichan, art contest winner on pixiv, 8K, intricate details, wet effects, rain drops, ethereal, mysterious, futuristic, UHD, HDR, cinematic lighting, in a beautiful forest, rainy day, award winning, trending on artstation, beautiful confident cheerful young woman, wearing a futuristic sleeveless dress, ultra beautiful detailed  eyes, hyper-detailed face, complex,  perfect, model,  textured,  chiaroscuro, professional make-up, realistic, figure in frame, "),
+        Image2ImageRegion(352-800, 352, 0, 352, reference_image=iic_image, strength=1.0),
+    ],
+    num_inference_steps=100,
+    seed=5525475061,
+)["images"][0]
+```
+
+![Input_Image](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/input_image.png)
+![mixture_canvas_results](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/canvas.png)
+
+### Stable Diffusion Mixture Tiling Pipeline SDXL
+
+This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
+
+```python
+import torch
+from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler, AutoencoderKL
+
+device="cuda"
+
+# Load fixed vae (optional)
+vae = AutoencoderKL.from_pretrained(
+    "madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16
+).to(device)
+
+# Create scheduler and model (similar to StableDiffusionPipeline)
+model_id="stablediffusionapi/yamermix-v8-vae"
+scheduler = DPMSolverMultistepScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
+pipe = DiffusionPipeline.from_pretrained(
+    model_id,
+    torch_dtype=torch.float16,
+    vae=vae,
+    custom_pipeline="mixture_tiling_sdxl",
+    scheduler=scheduler,
+    use_safetensors=False    
+).to(device)
+
+pipe.enable_model_cpu_offload()
+pipe.enable_vae_tiling()
+pipe.enable_vae_slicing()
+
+generator = torch.Generator(device).manual_seed(297984183)
+
+# Mixture of Diffusers generation
+image = pipe(
+    prompt=[[
+        "A charming house in the countryside, by jakub rozalski, sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece",
+        "A dirt road in the countryside crossing pastures, by jakub rozalski, sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece",        
+        "An old and rusty giant robot lying on a dirt road, by jakub rozalski, dark sunset lighting, elegant, highly detailed, smooth, sharp focus, artstation, stunning masterpiece"
+    ]],
+    tile_height=1024,
+    tile_width=1280,
+    tile_row_overlap=0,
+    tile_col_overlap=256,
+    guidance_scale_tiles=[[7, 7, 7]], # or guidance_scale=7 if is the same for all prompts
+    height=1024,
+    width=3840,    
+    generator=generator,
+    num_inference_steps=30,
+)["images"][0]
+```
+
+![mixture_tiling_results](https://huggingface.co/datasets/elismasilva/results/resolve/main/mixture_of_diffusers_sdxl_1.png)
+
 ### TensorRT Inpainting Stable Diffusion Pipeline
 
 The TensorRT Pipeline can be used to accelerate the Inpainting Stable Diffusion Inference run.
@@ -2457,41 +2672,6 @@ image = pipe(prompt, image=input_image, mask_image=mask_image, strength=0.75,).i
 image.save('tensorrt_inpaint_mecha_robot.png')
 ```
 
-### Stable Diffusion Mixture Canvas
-
-This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
-
-```python
-from PIL import Image
-from diffusers import LMSDiscreteScheduler, DiffusionPipeline
-from diffusers.pipelines.pipeline_utils import Image2ImageRegion, Text2ImageRegion, preprocess_image
-
-
-# Load and preprocess guide image
-iic_image = preprocess_image(Image.open("input_image.png").convert("RGB"))
-
-# Create scheduler and model (similar to StableDiffusionPipeline)
-scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
-pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler).to("cuda:0", custom_pipeline="mixture_canvas")
-pipeline.to("cuda")
-
-# Mixture of Diffusers generation
-output = pipeline(
-    canvas_height=800,
-    canvas_width=352,
-    regions=[
-        Text2ImageRegion(0, 800, 0, 352, guidance_scale=8,
-            prompt=f"best quality, masterpiece, WLOP, sakimichan, art contest winner on pixiv, 8K, intricate details, wet effects, rain drops, ethereal, mysterious, futuristic, UHD, HDR, cinematic lighting, in a beautiful forest, rainy day, award winning, trending on artstation, beautiful confident cheerful young woman, wearing a futuristic sleeveless dress, ultra beautiful detailed  eyes, hyper-detailed face, complex,  perfect, model,  textured,  chiaroscuro, professional make-up, realistic, figure in frame, "),
-        Image2ImageRegion(352-800, 352, 0, 352, reference_image=iic_image, strength=1.0),
-    ],
-    num_inference_steps=100,
-    seed=5525475061,
-)["images"][0]
-```
-
-![Input_Image](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/input_image.png)
-![mixture_canvas_results](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/canvas.png)
-
 ### IADB pipeline
 
 This pipeline is the implementation of the [α-(de)Blending: a Minimalist Deterministic Diffusion Model](https://arxiv.org/abs/2305.03486) paper.
@@ -3908,33 +4088,89 @@ This pipeline provides drag-and-drop image editing using stochastic differential
 See [paper](https://arxiv.org/abs/2311.01410), [paper page](https://ml-gsai.github.io/SDE-Drag-demo/), [original repo](https://github.com/ML-GSAI/SDE-Drag) for more information.
 
 ```py
-import PIL
 import torch
 from diffusers import DDIMScheduler, DiffusionPipeline
+from PIL import Image
+import requests
+from io import BytesIO
+import numpy as np
 
 # Load the pipeline
 model_path = "stable-diffusion-v1-5/stable-diffusion-v1-5"
 scheduler = DDIMScheduler.from_pretrained(model_path, subfolder="scheduler")
 pipe = DiffusionPipeline.from_pretrained(model_path, scheduler=scheduler, custom_pipeline="sde_drag")
-pipe.to('cuda')
 
-# To save GPU memory, torch.float16 can be used, but it may compromise image quality.
-# If not training LoRA, please avoid using torch.float16
-# pipe.to(torch.float16)
+# Ensure the model is moved to the GPU
+device = "cuda" if torch.cuda.is_available() else "cpu"
+pipe.to(device)
 
-# Provide prompt, image, mask image, and the starting and target points for drag editing.
-prompt = "prompt of the image"
-image = PIL.Image.open('/path/to/image')
-mask_image = PIL.Image.open('/path/to/mask_image')
-source_points = [[123, 456]]
-target_points = [[234, 567]]
+# Function to load image from URL
+def load_image_from_url(url):
+    response = requests.get(url)
+    return Image.open(BytesIO(response.content)).convert("RGB")
 
-# train_lora is optional, and in most cases, using train_lora can better preserve consistency with the original image.
-pipe.train_lora(prompt, image)
+# Function to prepare mask
+def prepare_mask(mask_image):
+    # Convert to grayscale
+    mask = mask_image.convert("L")
+    return mask
 
-output = pipe(prompt, image, mask_image, source_points, target_points)
-output_image = PIL.Image.fromarray(output)
+# Function to convert numpy array to PIL Image
+def array_to_pil(array):
+    # Ensure the array is in uint8 format
+    if array.dtype != np.uint8:
+        if array.max() <= 1.0:
+            array = (array * 255).astype(np.uint8)
+        else:
+            array = array.astype(np.uint8)
+    
+    # Handle different array shapes
+    if len(array.shape) == 3:
+        if array.shape[0] == 3:  # If channels first
+            array = array.transpose(1, 2, 0)
+        return Image.fromarray(array)
+    elif len(array.shape) == 4:  # If batch dimension
+        array = array[0]
+        if array.shape[0] == 3:  # If channels first
+            array = array.transpose(1, 2, 0)
+        return Image.fromarray(array)
+    else:
+        raise ValueError(f"Unexpected array shape: {array.shape}")
+
+# Image and mask URLs
+image_url = 'https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png'
+mask_url = 'https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png'
+
+# Load the images
+image = load_image_from_url(image_url)
+mask_image = load_image_from_url(mask_url)
+
+# Resize images to a size that's compatible with the model's latent space
+image = image.resize((512, 512))
+mask_image = mask_image.resize((512, 512))
+
+# Prepare the mask (keep as PIL Image)
+mask = prepare_mask(mask_image)
+
+# Provide the prompt and points for drag editing
+prompt = "A cute dog"
+source_points = [[32, 32]]  # Adjusted for 512x512 image
+target_points = [[64, 64]]  # Adjusted for 512x512 image
+
+# Generate the output image
+output_array = pipe(
+    prompt=prompt,
+    image=image,
+    mask_image=mask,
+    source_points=source_points,
+    target_points=target_points
+)
+
+# Convert output array to PIL Image and save
+output_image = array_to_pil(output_array)
 output_image.save("./output.png")
+print("Output image saved as './output.png'")
+
 ```
 
 ### Instaflow Pipeline
@@ -4585,8 +4821,8 @@ image = pipe(
 ```
 
 | ![Gradient](https://github.com/user-attachments/assets/e38ce4d5-1ae6-4df0-ab43-adc1b45716b5) | ![Input](https://github.com/user-attachments/assets/9c95679c-e9d7-4f5a-90d6-560203acd6b3) | ![Output](https://github.com/user-attachments/assets/5313ff64-a0c4-418b-8b55-a38f1a5e7532) |
-| ------------------------------------------------------------------------------------------ | --------------------------------------------------------------------------------------- | ---------------------------------------------------------------------------------------- |
-| Gradient                                                                                   | Input                                                                                   | Output                                                                                   |
+| -------------------------------------------------------------------------------------------- | ----------------------------------------------------------------------------------------- | ------------------------------------------------------------------------------------------ |
+| Gradient                                                                                     | Input                                                                                     | Output                                                                                     |
 
 A colab notebook demonstrating all results can be found [here](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing). Depth Maps have also been added in the same colab.
 
@@ -4634,6 +4870,93 @@ make_image_grid(image, rows=1, cols=len(image))
 # 50+, 100+, and 250+ num_inference_steps are recommended for nesting levels 0, 1, and 2 respectively.
 ```
 
+### Stable Diffusion XL Attentive Eraser Pipeline
+<img src="https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/fenmian.png"  width="600" />
+
+**Stable Diffusion XL Attentive Eraser Pipeline** is an advanced object removal pipeline that leverages SDXL for precise content suppression and seamless region completion. This pipeline uses **self-attention redirection guidance** to modify the model’s self-attention mechanism, allowing for effective removal and inpainting across various levels of mask precision, including semantic segmentation masks, bounding boxes, and hand-drawn masks. If you are interested in more detailed information and have any questions, please refer to the [paper](https://arxiv.org/abs/2412.12974) and [official implementation](https://github.com/Anonym0u3/AttentiveEraser).
+
+#### Key features
+
+- **Tuning-Free**: No additional training is required, making it easy to integrate and use.
+- **Flexible Mask Support**: Works with different types of masks for targeted object removal.
+- **High-Quality Results**: Utilizes the inherent generative power of diffusion models for realistic content completion.
+
+#### Usage example
+To use the Stable Diffusion XL Attentive Eraser Pipeline, you can initialize it as follows:
+```py
+import torch
+from diffusers import DDIMScheduler, DiffusionPipeline
+from diffusers.utils import load_image
+import torch.nn.functional as F
+from torchvision.transforms.functional import to_tensor, gaussian_blur
+
+dtype = torch.float16
+device = torch.device("cuda") if torch.cuda.is_available() else torch.device("cpu") 
+
+scheduler = DDIMScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False)
+pipeline = DiffusionPipeline.from_pretrained(
+    "stabilityai/stable-diffusion-xl-base-1.0",
+    custom_pipeline="pipeline_stable_diffusion_xl_attentive_eraser",
+    scheduler=scheduler,
+    variant="fp16",
+    use_safetensors=True,
+    torch_dtype=dtype,
+).to(device)
+
+
+def preprocess_image(image_path, device):
+    image = to_tensor((load_image(image_path)))
+    image = image.unsqueeze_(0).float() * 2 - 1 # [0,1] --> [-1,1]
+    if image.shape[1] != 3:
+        image = image.expand(-1, 3, -1, -1)
+        image = F.interpolate(image, (1024, 1024))
+        image = image.to(dtype).to(device)
+        return image
+
+def preprocess_mask(mask_path, device):
+    mask = to_tensor((load_image(mask_path, convert_method=lambda img: img.convert('L'))))
+    mask = mask.unsqueeze_(0).float()  # 0 or 1
+    mask = F.interpolate(mask, (1024, 1024))
+    mask = gaussian_blur(mask, kernel_size=(77, 77))
+    mask[mask < 0.1] = 0
+    mask[mask >= 0.1] = 1
+    mask = mask.to(dtype).to(device)
+    return mask
+
+prompt = "" # Set prompt to null
+seed=123 
+generator = torch.Generator(device=device).manual_seed(seed)
+source_image_path = "https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/an1024.png"
+mask_path = "https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/an1024_mask.png"
+source_image = preprocess_image(source_image_path, device)
+mask = preprocess_mask(mask_path, device)
+
+image = pipeline(
+    prompt=prompt, 
+    image=source_image,
+    mask_image=mask,
+    height=1024,
+    width=1024,
+    AAS=True, # enable AAS
+    strength=0.8, # inpainting strength
+    rm_guidance_scale=9, # removal guidance scale
+    ss_steps = 9, # similarity suppression steps
+    ss_scale = 0.3, # similarity suppression scale
+    AAS_start_step=0, # AAS start step
+    AAS_start_layer=34, # AAS start layer
+    AAS_end_layer=70, # AAS end layer
+    num_inference_steps=50, # number of inference steps # AAS_end_step = int(strength*num_inference_steps)
+    generator=generator,
+    guidance_scale=1,
+).images[0]
+image.save('./removed_img.png')
+print("Object removal completed")
+```
+
+| Source Image                                                                                   | Mask                                                                                        | Output                                                                                              |
+| ---------------------------------------------------------------------------------------------- | ------------------------------------------------------------------------------------------- | --------------------------------------------------------------------------------------------------- |
+| ![Source Image](https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/an1024.png) | ![Mask](https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/an1024_mask.png) | ![Output](https://raw.githubusercontent.com/Anonym0u3/Images/refs/heads/main/AE_step40_layer34.png) |
+
 # Perturbed-Attention Guidance
 
 [Project](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) / [arXiv](https://arxiv.org/abs/2403.17377) / [GitHub](https://github.com/KU-CVLAB/Perturbed-Attention-Guidance)

examples/community/fresco_v2v.py

@@ -404,10 +404,11 @@ def my_forward(
             # TODO: this requires sync between CPU and GPU. So try to pass timesteps as tensors if you can
             # This would be a good case for the `match` statement (Python 3.10+)
             is_mps = sample.device.type == "mps"
+            is_npu = sample.device.type == "npu"
             if isinstance(timestep, float):
-                dtype = torch.float32 if is_mps else torch.float64
+                dtype = torch.float32 if (is_mps or is_npu) else torch.float64
             else:
-                dtype = torch.int32 if is_mps else torch.int64
+                dtype = torch.int32 if (is_mps or is_npu) else torch.int64
             timesteps = torch.tensor([timesteps], dtype=dtype, device=sample.device)
         elif len(timesteps.shape) == 0:
             timesteps = timesteps[None].to(sample.device)

examples/community/matryoshka.py

@@ -80,7 +80,6 @@ from diffusers.utils import (
     USE_PEFT_BACKEND,
     BaseOutput,
     deprecate,
-    is_torch_version,
     is_torch_xla_available,
     logging,
     replace_example_docstring,
@@ -869,23 +868,7 @@ class CrossAttnDownBlock2D(nn.Module):
 
         for i, (resnet, attn) in enumerate(blocks):
             if torch.is_grad_enabled() and self.gradient_checkpointing:
-
-                def create_custom_forward(module, return_dict=None):
-                    def custom_forward(*inputs):
-                        if return_dict is not None:
-                            return module(*inputs, return_dict=return_dict)
-                        else:
-                            return module(*inputs)
-
-                    return custom_forward
-
-                ckpt_kwargs: Dict[str, Any] = {"use_reentrant": False} if is_torch_version(">=", "1.11.0") else {}
-                hidden_states = torch.utils.checkpoint.checkpoint(
-                    create_custom_forward(resnet),
-                    hidden_states,
-                    temb,
-                    **ckpt_kwargs,
-                )
+                hidden_states = self._gradient_checkpointing_func(resnet, hidden_states, temb)
                 hidden_states = attn(
                     hidden_states,
                     encoder_hidden_states=encoder_hidden_states,
@@ -1030,17 +1013,6 @@ class UNetMidBlock2DCrossAttn(nn.Module):
         hidden_states = self.resnets[0](hidden_states, temb)
         for attn, resnet in zip(self.attentions, self.resnets[1:]):
             if torch.is_grad_enabled() and self.gradient_checkpointing:
-
-                def create_custom_forward(module, return_dict=None):
-                    def custom_forward(*inputs):
-                        if return_dict is not None:
-                            return module(*inputs, return_dict=return_dict)
-                        else:
-                            return module(*inputs)
-
-                    return custom_forward
-
-                ckpt_kwargs: Dict[str, Any] = {"use_reentrant": False} if is_torch_version(">=", "1.11.0") else {}
                 hidden_states = attn(
                     hidden_states,
                     encoder_hidden_states=encoder_hidden_states,
@@ -1049,12 +1021,7 @@ class UNetMidBlock2DCrossAttn(nn.Module):
                     encoder_attention_mask=encoder_attention_mask,
                     return_dict=False,
                 )[0]
-                hidden_states = torch.utils.checkpoint.checkpoint(
-                    create_custom_forward(resnet),
-                    hidden_states,
-                    temb,
-                    **ckpt_kwargs,
-                )
+                hidden_states = self._gradient_checkpointing_func(resnet, hidden_states, temb)
             else:
                 hidden_states = attn(
                     hidden_states,
@@ -1192,23 +1159,7 @@ class CrossAttnUpBlock2D(nn.Module):
             hidden_states = torch.cat([hidden_states, res_hidden_states], dim=1)
 
             if torch.is_grad_enabled() and self.gradient_checkpointing:
-
-                def create_custom_forward(module, return_dict=None):
-                    def custom_forward(*inputs):
-                        if return_dict is not None:
-                            return module(*inputs, return_dict=return_dict)
-                        else:
-                            return module(*inputs)
-
-                    return custom_forward
-
-                ckpt_kwargs: Dict[str, Any] = {"use_reentrant": False} if is_torch_version(">=", "1.11.0") else {}
-                hidden_states = torch.utils.checkpoint.checkpoint(
-                    create_custom_forward(resnet),
-                    hidden_states,
-                    temb,
-                    **ckpt_kwargs,
-                )
+                hidden_states = self._gradient_checkpointing_func(resnet, hidden_states, temb)
                 hidden_states = attn(
                     hidden_states,
                     encoder_hidden_states=encoder_hidden_states,
@@ -1282,10 +1233,6 @@ class MatryoshkaTransformer2DModel(LegacyModelMixin, LegacyConfigMixin):
             ]
         )
 
-    def _set_gradient_checkpointing(self, module, value=False):
-        if hasattr(module, "gradient_checkpointing"):
-            module.gradient_checkpointing = value
-
     def forward(
         self,
         hidden_states: torch.Tensor,
@@ -1365,19 +1312,8 @@ class MatryoshkaTransformer2DModel(LegacyModelMixin, LegacyConfigMixin):
         # Blocks
         for block in self.transformer_blocks:
             if torch.is_grad_enabled() and self.gradient_checkpointing:
-
-                def create_custom_forward(module, return_dict=None):
-                    def custom_forward(*inputs):
-                        if return_dict is not None:
-                            return module(*inputs, return_dict=return_dict)
-                        else:
-                            return module(*inputs)
-
-                    return custom_forward
-
-                ckpt_kwargs: Dict[str, Any] = {"use_reentrant": False} if is_torch_version(">=", "1.11.0") else {}
-                hidden_states = torch.utils.checkpoint.checkpoint(
-                    create_custom_forward(block),
+                hidden_states = self._gradient_checkpointing_func(
+                    block,
                     hidden_states,
                     attention_mask,
                     encoder_hidden_states,
@@ -1385,7 +1321,6 @@ class MatryoshkaTransformer2DModel(LegacyModelMixin, LegacyConfigMixin):
                     timestep,
                     cross_attention_kwargs,
                     class_labels,
-                    **ckpt_kwargs,
                 )
             else:
                 hidden_states = block(
@@ -2724,10 +2659,6 @@ class MatryoshkaUNet2DConditionModel(
         for module in self.children():
             fn_recursive_set_attention_slice(module, reversed_slice_size)
 
-    def _set_gradient_checkpointing(self, module, value=False):
-        if hasattr(module, "gradient_checkpointing"):
-            module.gradient_checkpointing = value
-
     def enable_freeu(self, s1: float, s2: float, b1: float, b2: float):
         r"""Enables the FreeU mechanism from https://arxiv.org/abs/2309.11497.
 
@@ -2806,10 +2737,11 @@ class MatryoshkaUNet2DConditionModel(
             # TODO: this requires sync between CPU and GPU. So try to pass timesteps as tensors if you can
             # This would be a good case for the `match` statement (Python 3.10+)
             is_mps = sample.device.type == "mps"
+            is_npu = sample.device.type == "npu"
             if isinstance(timestep, float):
-                dtype = torch.float32 if is_mps else torch.float64
+                dtype = torch.float32 if (is_mps or is_npu) else torch.float64
             else:
-                dtype = torch.int32 if is_mps else torch.int64
+                dtype = torch.int32 if (is_mps or is_npu) else torch.int64
             timesteps = torch.tensor([timesteps], dtype=dtype, device=sample.device)
         elif len(timesteps.shape) == 0:
             timesteps = timesteps[None].to(sample.device)

examples/community/pipeline_flux_differential_img2img.py

@@ -87,7 +87,7 @@ def calculate_shift(
     base_seq_len: int = 256,
     max_seq_len: int = 4096,
     base_shift: float = 0.5,
-    max_shift: float = 1.16,
+    max_shift: float = 1.15,
 ):
     m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
     b = base_shift - m * base_seq_len
@@ -878,7 +878,7 @@ class FluxDifferentialImg2ImgPipeline(DiffusionPipeline, FluxLoraLoaderMixin):
             self.scheduler.config.get("base_image_seq_len", 256),
             self.scheduler.config.get("max_image_seq_len", 4096),
             self.scheduler.config.get("base_shift", 0.5),
-            self.scheduler.config.get("max_shift", 1.16),
+            self.scheduler.config.get("max_shift", 1.15),
         )
         timesteps, num_inference_steps = retrieve_timesteps(
             self.scheduler,

examples/community/pipeline_flux_rf_inversion.py

@@ -94,7 +94,7 @@ def calculate_shift(
     base_seq_len: int = 256,
     max_seq_len: int = 4096,
     base_shift: float = 0.5,
-    max_shift: float = 1.16,
+    max_shift: float = 1.15,
 ):
     m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
     b = base_shift - m * base_seq_len
@@ -823,7 +823,7 @@ class RFInversionFluxPipeline(
             self.scheduler.config.get("base_image_seq_len", 256),
             self.scheduler.config.get("max_image_seq_len", 4096),
             self.scheduler.config.get("base_shift", 0.5),
-            self.scheduler.config.get("max_shift", 1.16),
+            self.scheduler.config.get("max_shift", 1.15),
         )
         timesteps, num_inference_steps = retrieve_timesteps(
             self.scheduler,
@@ -993,7 +993,7 @@ class RFInversionFluxPipeline(
             self.scheduler.config.get("base_image_seq_len", 256),
             self.scheduler.config.get("max_image_seq_len", 4096),
             self.scheduler.config.get("base_shift", 0.5),
-            self.scheduler.config.get("max_shift", 1.16),
+            self.scheduler.config.get("max_shift", 1.15),
         )
         timesteps, num_inversion_steps = retrieve_timesteps(
             self.scheduler,

examples/community/pipeline_flux_with_cfg.py

@@ -70,7 +70,7 @@ def calculate_shift(
     base_seq_len: int = 256,
     max_seq_len: int = 4096,
     base_shift: float = 0.5,
-    max_shift: float = 1.16,
+    max_shift: float = 1.15,
 ):
     m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
     b = base_shift - m * base_seq_len
@@ -759,7 +759,7 @@ class FluxCFGPipeline(DiffusionPipeline, FluxLoraLoaderMixin, FromSingleFileMixi
             self.scheduler.config.get("base_image_seq_len", 256),
             self.scheduler.config.get("max_image_seq_len", 4096),
             self.scheduler.config.get("base_shift", 0.5),
-            self.scheduler.config.get("max_shift", 1.16),
+            self.scheduler.config.get("max_shift", 1.15),
         )
         timesteps, num_inference_steps = retrieve_timesteps(
             self.scheduler,

examples/community/stable_diffusion_tensorrt_img2img.py

@@ -1,5 +1,5 @@
 #
-# Copyright 2024 The HuggingFace Inc. team.
+# Copyright 2025 The HuggingFace Inc. team.
 # SPDX-FileCopyrightText: Copyright (c) 1993-2023 NVIDIA CORPORATION & AFFILIATES. All rights reserved.
 # SPDX-License-Identifier: Apache-2.0
 #

examples/community/stable_diffusion_tensorrt_inpaint.py

@@ -1,5 +1,5 @@
 #
-# Copyright 2024 The HuggingFace Inc. team.
+# Copyright 2025 The HuggingFace Inc. team.
 # SPDX-FileCopyrightText: Copyright (c) 1993-2023 NVIDIA CORPORATION & AFFILIATES. All rights reserved.
 # SPDX-License-Identifier: Apache-2.0
 #

examples/community/stable_diffusion_tensorrt_txt2img.py

@@ -1,5 +1,5 @@
 #
-# Copyright 2024 The HuggingFace Inc. team.
+# Copyright 2025 The HuggingFace Inc. team.
 # SPDX-FileCopyrightText: Copyright (c) 1993-2023 NVIDIA CORPORATION & AFFILIATES. All rights reserved.
 # SPDX-License-Identifier: Apache-2.0
 #

examples/community/stable_diffusion_xl_controlnet_reference.py

@@ -193,7 +193,8 @@ class StableDiffusionXLControlNetReferencePipeline(StableDiffusionXLControlNetPi
 
     def prepare_ref_latents(self, refimage, batch_size, dtype, device, generator, do_classifier_free_guidance):
         refimage = refimage.to(device=device)
-        if self.vae.dtype == torch.float16 and self.vae.config.force_upcast:
+        needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast
+        if needs_upcasting:
             self.upcast_vae()
             refimage = refimage.to(next(iter(self.vae.post_quant_conv.parameters())).dtype)
         if refimage.dtype != self.vae.dtype:
@@ -223,6 +224,11 @@ class StableDiffusionXLControlNetReferencePipeline(StableDiffusionXLControlNetPi
 
         # aligning device to prevent device errors when concating it with the latent model input
         ref_image_latents = ref_image_latents.to(device=device, dtype=dtype)
+
+        # cast back to fp16 if needed
+        if needs_upcasting:
+            self.vae.to(dtype=torch.float16)
+
         return ref_image_latents
 
     def prepare_ref_image(

examples/community/stable_diffusion_xl_reference.py

@@ -139,7 +139,8 @@ def retrieve_timesteps(
 class StableDiffusionXLReferencePipeline(StableDiffusionXLPipeline):
     def prepare_ref_latents(self, refimage, batch_size, dtype, device, generator, do_classifier_free_guidance):
         refimage = refimage.to(device=device)
-        if self.vae.dtype == torch.float16 and self.vae.config.force_upcast:
+        needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast
+        if needs_upcasting:
             self.upcast_vae()
             refimage = refimage.to(next(iter(self.vae.post_quant_conv.parameters())).dtype)
         if refimage.dtype != self.vae.dtype:
@@ -169,6 +170,11 @@ class StableDiffusionXLReferencePipeline(StableDiffusionXLPipeline):
 
         # aligning device to prevent device errors when concating it with the latent model input
         ref_image_latents = ref_image_latents.to(device=device, dtype=dtype)
+
+        # cast back to fp16 if needed
+        if needs_upcasting:
+            self.vae.to(dtype=torch.float16)
+
         return ref_image_latents
 
     def prepare_ref_image(

examples/consistency_distillation/train_lcm_distill_lora_sd_wds.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/consistency_distillation/train_lcm_distill_lora_sdxl_wds.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/consistency_distillation/train_lcm_distill_sd_wds.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/consistency_distillation/train_lcm_distill_sdxl_wds.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/controlnet/train_controlnet.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.
@@ -1143,7 +1143,7 @@ def main(args):
             if global_step >= args.max_train_steps:
                 break
 
-    # Create the pipeline using using the trained modules and save it.
+    # Create the pipeline using the trained modules and save it.
     accelerator.wait_for_everyone()
     if accelerator.is_main_process:
         controlnet = unwrap_model(controlnet)

examples/controlnet/train_controlnet_flax.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/controlnet/train_controlnet_flux.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.

examples/controlnet/train_controlnet_sd3.py

@@ -1,6 +1,6 @@
 #!/usr/bin/env python
 # coding=utf-8
-# Copyright 2024 The HuggingFace Inc. team. All rights reserved.
+# Copyright 2025 The HuggingFace Inc. team. All rights reserved.
 #
 # Licensed under the Apache License, Version 2.0 (the "License");
 # you may not use this file except in compliance with the License.