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139 Commits

Author SHA1 Message Date
yiyixuxu 78c125f76f up 2024-08-06 07:20:35 +02:00
yiyixuxu 9e3b0fe107 add dynamic cfg 2024-08-06 03:59:06 +02:00
Aryan 878f609aa5 remove dynamic guidance scale 2024-08-05 17:52:13 +02:00
zR 32da2e7673 Update lora_conversion_utils.py 2024-08-05 22:20:06 +08:00
zR 9a0b906518 restore 2024-08-05 22:08:50 +08:00
Aryan b3428ad5f5 Merge branch 'main' into cogvideox-2b 2024-08-05 18:09:09 +05:30
Aryan 70a54a8230 use num_frames instead of num_seconds 2024-08-05 14:38:42 +02:00
Aryan b7058d142c PAG variant for HunyuanDiT, PAG refactor (#8936)
* copy hunyuandit pipeline

* pag variant of hunyuan dit

* add tests

* update docs

* make style

* make fix-copies

* Update src/diffusers/pipelines/pag/pag_utils.py

* remove incorrect copied from

* remove pag hunyuan attn procs to resolve conflicts

* add pag attn procs again

* new implementation for pag_utils

* revert pag changes

* add pag refactor back; update pixart sigma

* update pixart pag tests

* apply suggestions from review

Co-Authored-By: yixu310@gmail.com

* make style

* update docs, fix tests

* fix tests

* fix test_components_function since list not accepted as valid __init__ param

* apply patch to fix broken tests

Co-Authored-By: Sayak Paul <spsayakpaul@gmail.com>

* make style

* fix hunyuan tests

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-05 17:56:09 +05:30
yiyixuxu 6f4e60b58b Merge branch 'cogvideox-2b' of github.com:zRzRzRzRzRzRzR/diffusers into cogvideo-dpm 2024-08-05 14:06:15 +02:00
yiyixuxu e4d65ccdd7 update 2024-08-05 13:57:42 +02:00
Vinh H. Pham e1d508ae92 [Tests] Improve transformers model test suite coverage - Latte (#8919)
* add LatteTransformer3DModel model test

* change patch_size to 1

* reduce req len

* reduce channel dims

* increase num_layers

* reduce dims further

* run make style

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Aryan <aryan@huggingface.co>
2024-08-05 17:13:03 +05:30
zR 22dcceb858 messages 2024-08-05 19:17:16 +08:00
zR 9c6b8894ff rename unsample and add some docs 2024-08-05 17:57:54 +08:00
zR cf7369d418 fix some error 2024-08-05 16:45:41 +08:00
Aryan 123ecef2b9 Merge branch 'main' into cogvideox-2b 2024-08-05 12:35:10 +05:30
zR 511c9ef560 merge remote branch 2024-08-05 14:02:36 +08:00
zR fb6130fe90 Merge branch 'cogvideox-common-draft-2' of github.com:huggingface/diffusers-new-model-addition-thudm into cogvideox-2b 2024-08-05 14:01:49 +08:00
zR 2d9602cc96 add CogVideoX team, Tsinghua University & ZhipuAI 2024-08-05 13:55:25 +08:00
Sayak Paul fc6a91e383 [FLUX] support LoRA (#9057)
* feat: lora support for Flux.

add tests

fix imports

major fixes.

* fix

fixes

final fixes?

* fix

* remove is_peft_available.
2024-08-05 10:24:05 +05:30
zR 1b1b737acb cogvideox branch 2024-08-05 11:28:31 +08:00
Aryan 311845fc77 update docs 2024-08-04 23:46:50 +02:00
Aryan 01c2dff338 update docs 2024-08-04 23:45:56 +02:00
Aryan 03580c07b9 remove cogvideox-specific attention processor 2024-08-04 23:37:31 +02:00
Aryan fd11c0fbee Merge branch 'main' into cogvideox-common-draft-2 2024-08-04 23:29:18 +02:00
Aryan 92c8c00756 make fix-copies 2024-08-04 23:29:11 +02:00
Aryan 5781e017dd Merge branch 'huggingface:main' into main 2024-08-05 02:22:36 +05:30
sayakpaul 90aa8be534 add workflow for rebasing with upstream automatically. 2024-08-04 22:50:51 +02:00
Sayak Paul 2f1b7870e2 Revert "add workflow to rebase with upstream main nightly." 2024-08-04 22:50:51 +02:00
sayakpaul 7360ea1d03 add upstream 2024-08-04 22:50:51 +02:00
sayakpaul ba1855c07e add workflow to rebase with upstream main nightly. 2024-08-04 22:50:51 +02:00
Aryan 2b76099610 [refactor] apply qk norm in attention processors (#9071)
* apply qk norm in attention processors

* revert attention processor

* qk-norm in only attention proc 2.0 and fused variant
2024-08-04 05:42:46 -10:00
Aryan 1b1b26b65c make style 2024-08-04 16:29:57 +02:00
Aryan 312f7dc4fd apply suggestions from review 2024-08-04 16:29:35 +02:00
psychedelicious 4f0d01d387 type get_attention_scores as optional in get_attention_scores (#9075)
`None` is valid for `get_attention_scores`, should be typed as such
2024-08-04 17:19:05 +05:30
asfiyab-nvidia 3dc10a535f Update TensorRT txt2img and inpaint community pipelines (#9037)
* Update TensorRT txt2img and inpaint community pipelines

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

* update tensorrt install instructions

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

---------

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-04 16:00:40 +05:30
Sayak Paul c370b90ff1 [Flux] minor documentation fixes for flux. (#9048)
* minor documentation fixes for flux.

* clipskip

* add gist
2024-08-04 15:53:01 +05:30
zR ba4223ac3b remove attenstion mask 2024-08-04 15:52:37 +08:00
zR 6988cc3a86 new schedule with dpm 2024-08-04 14:26:00 +08:00
Philip Rideout ebf3ab1477 Fix grammar mistake. (#9072) 2024-08-04 04:32:03 +05:30
Aryan fa7fa9cced make inference 2-3x faster (by fixing the bug i introduced) 🚀😎 2024-08-03 18:48:54 +02:00
Aryan 61c6da076a remove chunked ff code; reuse and refactor to support temb directly in adalayernorm 2024-08-03 17:43:09 +02:00
zR c7ee165c4f Update downsampling.py 2024-08-03 22:22:28 +08:00
zR d99528be94 Restore the timesteps parameter 2024-08-03 16:04:46 +08:00
Aryan fbe29c6298 [refactor] create modeling blocks specific to AnimateDiff (#8979)
* animatediff specific transformer model

* make style

* make fix-copies

* move blocks to unet motion model

* make style

* remove dummy object

* fix incorrectly passed param causing test failures

* rename model and output class

* fix sparsectrl imports

* remove todo comments

* remove temporal double self attn param from controlnet sparsectrl

* add deprecated versions of blocks

* apply suggestions from review

* update

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-08-03 13:03:39 +05:30
zR fd0831c52c timestep fix 2024-08-03 15:11:12 +08:00
zR 477e12b235 fix 2024-08-03 15:08:07 +08:00
zR b42b079213 fix some comment 2024-08-03 13:46:57 +08:00
Tolga Cangöz 7071b7461b Errata: Fix typos & \s+$ (#9008)
* Fix typos

* chore: Fix typos

* chore: Update README.md for promptdiffusion example

* Trim trailing white spaces

* Fix a typo

* update number

* chore: update number

* Trim trailing white space

* Update README.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update README.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

---------

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2024-08-02 21:24:25 -07:00
Frank (Haofan) Wang a054c78495 Update transformer_flux.py (#9060) 2024-08-03 08:58:32 +05:30
zR 21509aa7f5 fp16 problem 2024-08-03 00:34:55 +08:00
Aryan 65f6211f1f Merge pull request #4 from huggingface/cogvideox-refactor-to-diffusers
CogVideoX pipeline, refactor modeling to diffusers format, bug fixes
2024-08-02 21:20:34 +05:30
Aryan 3def90523d update 2024-08-02 17:40:36 +02:00
Aryan 551c884acd remove debug prints 2024-08-02 17:40:23 +02:00
zR ec53a30a0e schedule 2024-08-02 22:38:58 +08:00
Dhruv Nair b1f43d7189 Fix Nightly Deps (#9036)
update
2024-08-02 15:18:54 +05:30
zR 71e7c82ae8 vae problem fix 2024-08-02 16:23:24 +08:00
Sayak Paul 0e460675e2 [Flux] allow tests to run (#9050)
* fix tests

* fix

* float64 skip

* remove sample_size.

* remove

* remove more

* default_sample_size.

* credit black forest for flux model.

* skip

* fix: tests

* remove OriginalModelMixin

* add transformer model test

* add: transformer model tests
2024-08-02 11:49:59 +05:30
Sayak Paul 7b98c4cc67 [Core] Add PAG support for PixArtSigma (#8921)
* feat: add pixart sigma pag.

* inits.

* fixes

* fix

* remove print.

* copy paste methods to the pixart pag mixin

* fix-copies

* add documentation.

* add tests.

* remove correction file.

* remove pag_applied_layers

* empty
2024-08-02 07:12:41 +05:30
Aryan c33dd0213b remove incorrect copied from 2024-08-02 01:04:35 +02:00
Aryan e12458e16c make style 2024-08-02 01:03:37 +02:00
Aryan 77558f31bf add pipeline docs 2024-08-02 01:03:02 +02:00
Aryan 41da084fbe remove debug prints 2024-08-02 00:53:09 +02:00
Aryan 4c2e8870e6 add cogvideo specific attn processor 2024-08-02 00:51:15 +02:00
Sayak Paul 27637a5402 Flux pipeline (#9043)
add flux!

Signed-off-by: Adrien <adrien@huggingface.co>
Co-authored-by: Adrien <adrien.69740@gmail.com>
Co-authored-by: Anatoly Belikov <abelikov@singularitynet.io>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: yiyixuxu <yixu310@gmail.com>
2024-08-01 11:30:52 -10:00
Aryan 2ea22e1cc7 [docs] fix pia example (#9015)
fix pia example docstring
2024-08-02 02:47:40 +05:30
Aryan fe6f5d6419 ensure tokenizer config correctly uses 226 as text length 2024-08-01 22:06:08 +02:00
Aryan d0b8db2b11 update 2024-08-01 21:12:49 +02:00
zR 351d1f009e remove 0.transformer_blocks.encoder.embed_tokens.weight 2024-08-01 23:31:03 +08:00
zR a31db5f952 using with 226 instead of 225 of final weight 2024-08-01 22:58:27 +08:00
YiYi Xu 95a7832879 fix load sharded checkpoint from a subfolder (local path) (#8913)
fix

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-01 20:15:42 +05:30
Aryan 03ee7cd109 add pipeline implementation 2024-08-01 13:52:58 +02:00
Aryan 712ddbeac6 make style 2024-08-01 12:18:18 +02:00
Aryan 03c28eef5b modeling fixes 2024-08-01 12:18:01 +02:00
Dhruv Nair c646fbc124 Updates deps for pipeline test fetcher (#9033)
update
2024-08-01 13:22:13 +05:30
Aryan 05b706c003 PAG variant for AnimateDiff (#8789)
* add animatediff pag pipeline

* remove unnecessary print

* make fix-copies

* fix ip-adapter bug

* update docs

* add fast tests and fix bugs

* update

* update

* address review comments

* update ip adapter single test expected slice

* implement test_from_pipe_consistent_config; fix expected slice values

* LoraLoaderMixin->StableDiffusionLoraLoaderMixin; add latest freeinit test
2024-08-01 12:39:39 +05:30
Aryan e05f83479c refactor vae 2024-08-01 06:05:49 +02:00
zR bb4740ce29 add clear_fake_cp_cache 2024-08-01 00:26:59 +08:00
zR 2956866ef4 Merge branch 'cogvideox-common-draft-2' of https://github.com/huggingface/diffusers-new-model-addition-thudm into cogvideox-common-draft-2 2024-07-31 21:41:16 +08:00
zR 4498cfc98c add doc draft 2024-07-31 21:41:09 +08:00
Aryan a449ceb3ef update conversion script 2024-07-31 14:40:16 +02:00
Aryan 45f7127ade fix bug in handling long prompts 2024-07-31 14:28:46 +02:00
Aryan 73469f9562 make fix-copies 2024-07-31 14:13:56 +02:00
Aryan d45d199b99 make style 2024-07-31 14:13:38 +02:00
Aryan e67cc5ae47 implement encode prompt 2024-07-31 14:12:22 +02:00
Aryan 470815cefa minor refactor 2024-07-31 13:38:23 +02:00
Aryan 5f183bfe27 reorder upsampling/downsampling blocks in order of invocation 2024-07-31 13:33:57 +02:00
Aryan c43a8f5b2b minor factor and repositioning of code in order of invocation 2024-07-31 13:27:15 +02:00
Aryan 9f9d0cbb83 verify CogVideoXSpatialNorm3D implementation 2024-07-31 13:14:39 +02:00
Aryan 2be7469821 groups->norm_num_groups 2024-07-31 12:44:18 +02:00
Aryan 3ae9413966 undo unnecessary changes added on cogvideo-vae by mistake 2024-07-31 12:42:03 +02:00
Aryan ec9508c83b Merge branch 'main' into cogvideox-common-draft-2 2024-07-31 12:40:26 +02:00
Aryan 6bcafcbaa6 make fix-copies 2024-07-31 12:02:58 +02:00
Aryan b3052807e5 add skeleton for pipeline 2024-07-31 12:02:38 +02:00
Aryan 73b041e7a9 Merge branch 'cogvideox' into cogvideox-common-draft-1 2024-07-31 11:41:28 +02:00
zR 1c661ce3d4 remove deriving and using nn.module 2024-07-31 17:19:22 +08:00
zR 8fe54bcd26 add 2024-07-31 16:42:47 +08:00
zR ee40f0e1ca follow review guide 2024-07-31 14:31:21 +08:00
sayakpaul 0980f4dcd2 add workflow for rebasing with upstream automatically. 2024-07-31 08:35:22 +05:30
Sayak Paul 71bcb1e1c5 Revert "add workflow to rebase with upstream main nightly." 2024-07-31 08:35:22 +05:30
sayakpaul dfeb32975d add upstream 2024-07-31 08:35:22 +05:30
sayakpaul d83c1f8447 add workflow to rebase with upstream main nightly. 2024-07-31 08:35:22 +05:30
Aryan 21a0fc1b0d make style 2024-07-31 02:55:08 +02:00
Aryan 16967589d8 remove debug prints 2024-07-31 02:54:41 +02:00
Aryan d963b1aaa4 update conversion script for latest modeling changes 2024-07-31 02:54:32 +02:00
Aryan e982881716 refactor 2024-07-31 02:52:53 +02:00
Aryan cb5348a0c2 fix nasty bug in 3d sincos pos embeds 2024-07-31 01:26:32 +02:00
Yoach Lacombe ea1b4ea7ca Fix Stable Audio repository id (#9016)
Fix Stable Audio repo id
2024-07-30 23:17:44 +05:30
zR aff72ec5dc Update autoencoder_kl3d.py 2024-07-30 22:43:09 +08:00
zR dc7e6e814f fix error 2024-07-30 22:33:34 +08:00
zR a3d827fb8d rename 2024-07-30 22:00:53 +08:00
zR 84ff56eb90 fix with some review guide 2024-07-30 21:05:04 +08:00
Aryan 45cb1f92d3 fix layernorms 2024-07-30 14:57:57 +02:00
Aryan e5b94b4c57 [core] Move community AnimateDiff ControlNet to core (#8972)
* add animatediff controlnet to core

* make style; remove unused method

* fix copied from comment

* add tests

* changes to make tests work

* add utility function to load videos

* update docs

* update pipeline example

* make style

* update docs with example

* address review comments

* add latest freeinit test from #8969

* LoraLoaderMixin -> StableDiffusionLoraLoaderMixin

* fix docs

* Update src/diffusers/utils/loading_utils.py

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>

* fix: variable out of scope

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-30 17:10:37 +05:30
Yoach Lacombe 69e72b1dd1 Stable Audio integration (#8716)
* WIP modeling code and pipeline

* add custom attention processor + custom activation + add to init

* correct ProjectionModel forward

* add stable audio to __initèè

* add autoencoder and update pipeline and modeling code

* add half Rope

* add partial rotary v2

* add temporary modfis to scheduler

* add EDM DPM Solver

* remove TODOs

* clean GLU

* remove att.group_norm to attn processor

* revert back src/diffusers/schedulers/scheduling_dpmsolver_multistep.py

* refactor GLU -> SwiGLU

* remove redundant args

* add channel multiples in autoencoder docstrings

* changes in docsrtings and copyright headers

* clean pipeline

* further cleaning

* remove peft and lora and fromoriginalmodel

* Delete src/diffusers/pipelines/stable_audio/diffusers.code-workspace

* make style

* dummy models

* fix copied from

* add fast oobleck tests

* add brownian tree

* oobleck autoencoder slow tests

* remove TODO

* fast stable audio pipeline tests

* add slow tests

* make style

* add first version of docs

* wrap is_torchsde_available to the scheduler

* fix slow test

* test with input waveform

* add input waveform

* remove some todos

* create stableaudio gaussian projection + make style

* add pipeline to toctree

* fix copied from

* make quality

* refactor timestep_features->time_proj

* refactor joint_attention_kwargs->cross_attention_kwargs

* remove forward_chunk

* move StableAudioDitModel to transformers folder

* correct convert + remove partial rotary embed

* apply suggestions from yiyixuxu -> removing attn.kv_heads

* remove temb

* remove cross_attention_kwargs

* further removal of cross_attention_kwargs

* remove text encoder autocast to fp16

* continue removing autocast

* make style

* refactor how text and audio are embedded

* add paper

* update example code

* make style

* unify projection model forward + fix device placement

* make style

* remove fuse qkv

* apply suggestions from review

* Update src/diffusers/pipelines/stable_audio/pipeline_stable_audio.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* make style

* smaller models in fast tests

* pass sequential offloading fast tests

* add docs for vae and autoencoder

* make style and update example

* remove useless import

* add cosine scheduler

* dummy classes

* cosine scheduler docs

* better description of scheduler

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
2024-07-30 15:29:06 +05:30
Aryan 59e6669f6d fix attention mask 2024-07-30 11:44:42 +02:00
Aryan bb917755ee add imports 2024-07-30 09:04:45 +02:00
Aryan bd6efd5fe4 initial draft of cogvideo transformer 2024-07-30 09:02:06 +02:00
zR c341786f3e vae draft 2024-07-30 13:25:08 +08:00
zR c8e5491be0 Create autoencoder_kl3d.py 2024-07-30 13:15:58 +08:00
Sayak Paul 8c4856cd6c [LoRA] fix: animate diff lora stuff. (#8995)
* fix: animate diff lora stuff.

* fix scaling function for UNetMotionModel

* emoty
2024-07-30 09:18:41 +05:30
Anatoly Belikov f240a936da handle lora scale and clip skip in lpw sd and sdxl community pipelines (#8988)
* handle lora scale and clip skip in lpw sd and sdxl

* use StableDiffusionLoraLoaderMixin

* use StableDiffusionXLLoraLoaderMixin

* style

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-07-30 07:19:28 +05:30
Sayak Paul 00d8d46e23 [Docs] credit where it's due for Lumina and Latte. (#9000)
credit where it's due for Lumina and Latte.
2024-07-29 10:02:03 -07:00
Adrien bfc9369f0a [CI] Update runner configuration for setup and nightly tests (#9005)
* [CI] Update runner configuration for setup and nightly tests

Signed-off-by: Adrien <adrien@huggingface.co>

* fix group

Signed-off-by: Adrien <adrien@huggingface.co>

* update for t4

Signed-off-by: Adrien <adrien@huggingface.co>

---------

Signed-off-by: Adrien <adrien@huggingface.co>
2024-07-29 21:14:31 +05:30
Álvaro Somoza 73acebb8cf [Kolors] Add IP Adapter (#8901)
* initial draft

* apply suggestions

* fix failing test

* added ipa to img2img

* add docs

* apply suggestions
2024-07-26 14:25:44 -04:00
Aryan ca0747a07e remove unused code from pag attn procs (#8928) 2024-07-26 07:58:40 -10:00
Aryan 5c53ca5ed8 [core] AnimateDiff SparseCtrl (#8897)
* initial sparse control model draft

* remove unnecessary implementation

* copy animatediff pipeline

* remove deprecated callbacks

* update

* update pipeline implementation progress

* make style

* make fix-copies

* update progress

* add partially working pipeline

* remove debug prints

* add model docs

* dummy objects

* improve motion lora conversion script

* fix bugs

* update docstrings

* remove unnecessary model params; docs

* address review comment

* add copied from to zero_module

* copy animatediff test

* add fast tests

* update docs

* update

* update pipeline docs

* fix expected slice values

* fix license

* remove get_down_block usage

* remove temporal_double_self_attention from get_down_block

* update

* update docs with org and documentation images

* make from_unet work in sparsecontrolnetmodel

* add latest freeinit test from #8969

* make fix-copies

* LoraLoaderMixin -> StableDiffsuionLoraLoaderMixin
2024-07-26 17:46:05 +05:30
Aryan 57a021d5e4 [fix] FreeInit step index out of bounds (#8969)
* fix step index out of bounds

* add test for free_init with different schedulers

* add test to vid2vid and pia
2024-07-26 16:45:55 +05:30
Dhruv Nair 1168eaaadd [CI] Nightly Test Runner explicitly set runner for Setup Pipeline Matrix (#8986)
* update

* update

* update
2024-07-26 13:20:35 +05:30
Dhruv Nair bce9105ac7 [CI] Fix parallelism in nightly tests (#8983)
update
2024-07-26 10:04:01 +05:30
RandomGamingDev 2afb2e0aac Added accelerator based gradient accumulation for basic_example (#8966)
added accelerator based gradient accumulation for basic_example

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-07-26 09:35:52 +05:30
Sayak Paul d87fe95f90 [Chore] add LoraLoaderMixin to the inits (#8981)
* introduce  to promote reusability.

* up

* add more tests

* up

* remove comments.

* fix fuse_nan test

* clarify the scope of fuse_lora and unfuse_lora

* remove space

* rewrite fuse_lora a bit.

* feedback

* copy over load_lora_into_text_encoder.

* address dhruv's feedback.

* fix-copies

* fix issubclass.

* num_fused_loras

* fix

* fix

* remove mapping

* up

* fix

* style

* fix-copies

* change to SD3TransformerLoRALoadersMixin

* Apply suggestions from code review

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>

* up

* handle wuerstchen

* up

* move lora to lora_pipeline.py

* up

* fix-copies

* fix documentation.

* comment set_adapters().

* fix-copies

* fix set_adapters() at the model level.

* fix?

* fix

* loraloadermixin.

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-26 08:59:33 +05:30
Sayak Paul 50e66f2f95 [Chore] remove all is from auraflow. (#8980)
remove all is from auraflow.
2024-07-26 07:31:06 +05:30
efwfe 9b8c8605d1 fix guidance_scale value not equal to the value in comments (#8941)
fix guidance_scale value not equal with the value in comments
2024-07-25 12:31:37 -10:00
YiYi Xu 62863bb1ea Revert "[LoRA] introduce LoraBaseMixin to promote reusability." (#8976)
Revert "[LoRA] introduce LoraBaseMixin to promote reusability. (#8774)"

This reverts commit 527430d0a4.
2024-07-25 09:10:35 -10:00
mazharosama 1fd647f2a0 Enable CivitAI SDXL Inpainting Models Conversion (#8795)
modify in_channels in network_config params

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-25 07:44:57 -10:00
asfiyab-nvidia 0bda1d7b89 Update TensorRT img2img community pipeline (#8899)
* Update TensorRT img2img pipeline

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

* Update TensorRT version installed

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

* make style and quality

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

* Update examples/community/stable_diffusion_tensorrt_img2img.py

Co-authored-by: Tolga Cangöz <46008593+tolgacangoz@users.noreply.github.com>

* Update examples/community/README.md

Co-authored-by: Tolga Cangöz <46008593+tolgacangoz@users.noreply.github.com>

* Apply style and quality using ruff 0.1.5

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>

---------

Signed-off-by: Asfiya Baig <asfiyab@nvidia.com>
Co-authored-by: Tolga Cangöz <46008593+tolgacangoz@users.noreply.github.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-25 21:58:21 +05:30
Sayak Paul 527430d0a4 [LoRA] introduce LoraBaseMixin to promote reusability. (#8774)
* introduce  to promote reusability.

* up

* add more tests

* up

* remove comments.

* fix fuse_nan test

* clarify the scope of fuse_lora and unfuse_lora

* remove space

* rewrite fuse_lora a bit.

* feedback

* copy over load_lora_into_text_encoder.

* address dhruv's feedback.

* fix-copies

* fix issubclass.

* num_fused_loras

* fix

* fix

* remove mapping

* up

* fix

* style

* fix-copies

* change to SD3TransformerLoRALoadersMixin

* Apply suggestions from code review

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>

* up

* handle wuerstchen

* up

* move lora to lora_pipeline.py

* up

* fix-copies

* fix documentation.

* comment set_adapters().

* fix-copies

* fix set_adapters() at the model level.

* fix?

* fix

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-25 21:40:58 +05:30
Aryan 3ae0ee88d3 [tests] speed up animatediff tests (#8846)
* speed up animatediff tests

* fix pia test_ip_adapter_single

* fix tests/pipelines/pia/test_pia.py::PIAPipelineFastTests::test_dict_tuple_outputs_equivalent

* update

* fix ip adapter tests

* skip test_from_pipe_consistent_config tests

* fix prompt_embeds test

* update test_from_pipe_consistent_config tests

* fix expected_slice values

* remove temporal_norm_num_groups from UpBlockMotion

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2024-07-25 17:35:43 +05:30
Dhruv Nair 5fbb4d32d5 [CI] Slow Test Updates (#8870)
* update

* update

* update
2024-07-25 16:00:43 +05:30
259 changed files with 26498 additions and 5076 deletions
+3 -3
View File
@@ -13,13 +13,13 @@ env:
jobs:
torch_pipelines_cuda_benchmark_tests:
env:
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL_BENCHMARK }}
name: Torch Core Pipelines CUDA Benchmarking Tests
strategy:
fail-fast: false
max-parallel: 1
runs-on:
runs-on:
group: aws-g6-4xlarge-plus
container:
image: diffusers/diffusers-pytorch-compile-cuda
@@ -59,7 +59,7 @@ jobs:
if: ${{ success() }}
run: |
pip install requests && python utils/notify_benchmarking_status.py --status=success
- name: Report failure status
if: ${{ failure() }}
run: |
+5 -3
View File
@@ -20,7 +20,8 @@ env:
jobs:
test-build-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
if: github.event_name == 'pull_request'
steps:
- name: Set up Docker Buildx
@@ -50,7 +51,8 @@ jobs:
if: steps.file_changes.outputs.all != ''
build-and-push-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
if: github.event_name != 'pull_request'
permissions:
@@ -98,4 +100,4 @@ jobs:
slack_channel: ${{ env.CI_SLACK_CHANNEL }}
title: "🤗 Results of the ${{ matrix.image-name }} Docker Image build"
status: ${{ job.status }}
slack_token: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
slack_token: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
@@ -24,7 +24,7 @@ jobs:
mirror_community_pipeline:
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL_COMMUNITY_MIRROR }}
runs-on: ubuntu-latest
steps:
# Checkout to correct ref
@@ -95,7 +95,7 @@ jobs:
if: ${{ success() }}
run: |
pip install requests && python utils/notify_community_pipelines_mirror.py --status=success
- name: Report failure status
if: ${{ failure() }}
run: |
+29 -90
View File
@@ -7,7 +7,7 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
HF_HUB_ENABLE_HF_TRANSFER: 1
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
@@ -18,8 +18,11 @@ env:
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines Matrix
runs-on: diffusers/diffusers-pytorch-cpu
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
@@ -27,13 +30,9 @@ jobs:
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
pip install -e .
pip install -e .[test]
pip install huggingface_hub
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
@@ -50,16 +49,18 @@ jobs:
path: reports
run_nightly_tests_for_torch_pipelines:
name: Torch Pipelines CUDA Nightly Tests
name: Nightly Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -67,19 +68,16 @@ jobs:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: |
python utils/print_env.py
- name: Nightly PyTorch CUDA checkpoint (pipelines) tests
- name: Pipeline CUDA Test
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
@@ -90,38 +88,37 @@ jobs:
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
--report-log=tests_pipeline_${{ matrix.module }}_cuda.log \
tests/pipelines/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_for_other_torch_modules:
name: Torch Non-Pipelines CUDA Nightly Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
name: Nightly Torch CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
strategy:
max-parallel: 2
matrix:
module: [models, schedulers, others, examples]
module: [models, schedulers, lora, others, single_file, examples]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -133,8 +130,8 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
@@ -158,7 +155,6 @@ jobs:
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v --make-reports=examples_torch_cuda \
--report-log=examples_torch_cuda.log \
@@ -181,64 +177,7 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_lora_nightly_tests:
name: Nightly LoRA Tests with PEFT and TORCH
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly LoRA tests with PEFT and Torch
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_lora_cuda \
--report-log=tests_torch_lora_cuda.log \
tests/lora
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_lora_cuda_stats.txt
cat reports/tests_torch_lora_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_lora_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_flax_tpu_tests:
name: Nightly Flax TPU Tests
@@ -294,14 +233,15 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_onnx_tests:
name: Nightly ONNXRuntime CUDA tests on Ubuntu
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
@@ -318,11 +258,10 @@ jobs:
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly ONNXRuntime CUDA tests
- name: Run Nightly ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
@@ -349,7 +288,7 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
@@ -411,4 +350,4 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
+7 -4
View File
@@ -15,7 +15,8 @@ concurrency:
jobs:
setup_pr_tests:
name: Setup PR Tests
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -73,7 +74,8 @@ jobs:
max-parallel: 2
matrix:
modules: ${{ fromJson(needs.setup_pr_tests.outputs.matrix) }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -123,12 +125,13 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
runs-on:
group: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
+3 -2
View File
@@ -71,7 +71,8 @@ jobs:
name: LoRA - ${{ matrix.lib-versions }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
@@ -128,4 +129,4 @@ jobs:
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports
path: reports
+8 -6
View File
@@ -77,28 +77,29 @@ jobs:
config:
- name: Fast PyTorch Pipeline CPU tests
framework: pytorch_pipelines
runner: [ self-hosted, intel-cpu, 32-cpu, 256-ram, ci ]
runner: aws-highmemory-32-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_pipelines
- name: Fast PyTorch Models & Schedulers CPU tests
framework: pytorch_models
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_models_schedulers
- name: Fast Flax CPU tests
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: PyTorch Example CPU tests
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
runs-on:
group: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
@@ -180,7 +181,8 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner:
group: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
+22 -71
View File
@@ -11,17 +11,16 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
PIPELINE_USAGE_CUTOFF: 50000
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
@@ -52,17 +51,18 @@ jobs:
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Slow Tests
name: Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -103,10 +103,11 @@ jobs:
torch_cuda_tests:
name: Torch CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
@@ -124,12 +125,13 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow PyTorch CUDA tests
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
@@ -153,61 +155,6 @@ jobs:
name: torch_cuda_test_reports
path: reports
peft_cuda_tests:
name: PEFT CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow PEFT CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda \
tests/lora/
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "lora and not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda_models_lora \
tests/models/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_peft_cuda_stats.txt
cat reports/tests_peft_cuda_failures_short.txt
cat reports/tests_peft_cuda_models_lora_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_peft_test_reports
path: reports
flax_tpu_tests:
name: Flax TPU Tests
runs-on: docker-tpu
@@ -257,7 +204,8 @@ jobs:
onnx_cuda_tests:
name: ONNX CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --gpus 0
@@ -305,11 +253,12 @@ jobs:
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
@@ -347,11 +296,12 @@ jobs:
run_xformers_tests:
name: PyTorch xformers CUDA tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-xformers-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
@@ -388,11 +338,12 @@ jobs:
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
+6 -5
View File
@@ -29,28 +29,29 @@ jobs:
config:
- name: Fast PyTorch CPU tests on Ubuntu
framework: pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu
- name: Fast Flax CPU tests on Ubuntu
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: Fast ONNXRuntime CPU tests on Ubuntu
framework: onnxruntime
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-onnxruntime-cpu
report: onnx_cpu
- name: PyTorch Example CPU tests on Ubuntu
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: aws-general-8-plus
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
runs-on:
group: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
+3 -2
View File
@@ -26,7 +26,8 @@ env:
jobs:
run_tests:
name: "Run a test on our runner from a PR"
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on:
group: aws-g4dn-2xlarge
container:
image: ${{ github.event.inputs.docker_image }}
options: --gpus 0 --privileged --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
@@ -70,4 +71,4 @@ jobs:
env:
PY_TEST: ${{ github.event.inputs.test }}
run: |
pytest "$PY_TEST"
pytest "$PY_TEST"
+2 -1
View File
@@ -19,7 +19,8 @@ env:
jobs:
ssh_runner:
name: "SSH"
runs-on: [self-hosted, intel-cpu, 32-cpu, 256-ram, ci]
runs-on:
group: aws-highmemory-32-plus
container:
image: ${{ github.event.inputs.docker_image }}
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --privileged
+2 -1
View File
@@ -22,7 +22,8 @@ env:
jobs:
ssh_runner:
name: "SSH"
runs-on: [single-gpu, nvidia-gpu, "${{ github.event.inputs.runner_type }}", ci]
runs-on:
group: "${{ github.event.inputs.runner_type }}"
container:
image: ${{ github.event.inputs.docker_image }}
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0 --privileged
+1 -1
View File
@@ -63,7 +63,7 @@ In the same spirit, you are of immense help to the community by answering such q
**Please** keep in mind that the more effort you put into asking or answering a question, the higher
the quality of the publicly documented knowledge. In the same way, well-posed and well-answered questions create a high-quality knowledge database accessible to everybody, while badly posed questions or answers reduce the overall quality of the public knowledge database.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formatted/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
**NOTE about channels**:
[*The forum*](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) is much better indexed by search engines, such as Google. Posts are ranked by popularity rather than chronologically. Hence, it's easier to look up questions and answers that we posted some time ago.
+2 -2
View File
@@ -67,7 +67,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 27.000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 30,000+ checkpoints):
```python
from diffusers import DiffusionPipeline
@@ -209,7 +209,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +12.000 other amazing GitHub repositories 💪
- +14,000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
@@ -38,6 +38,7 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
@@ -38,6 +38,7 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
+1
View File
@@ -38,6 +38,7 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
@@ -38,6 +38,7 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
+14
View File
@@ -239,6 +239,8 @@
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/autoencoder_oobleck
title: Oobleck AutoEncoder
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
@@ -251,6 +253,8 @@
title: HunyuanDiT2DModel
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
@@ -259,6 +263,8 @@
title: TransformerTemporalModel
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/stable_audio_transformer
title: StableAudioDiTModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/controlnet
@@ -267,6 +273,8 @@
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
title: Models
- isExpanded: false
sections:
@@ -314,6 +322,8 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/i2vgenxl
@@ -360,6 +370,8 @@
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/stable_audio
title: Stable Audio
- local: api/pipelines/stable_cascade
title: Stable Cascade
- sections:
@@ -423,6 +435,8 @@
title: CMStochasticIterativeScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/cosine_dpm
title: CosineDPMSolverMultistepScheduler
- local: api/schedulers/ddim_inverse
title: DDIMInverseScheduler
- local: api/schedulers/ddim
+21 -6
View File
@@ -12,10 +12,13 @@ specific language governing permissions and limitations under the License.
# LoRA
LoRA is a fast and lightweight training method that inserts and trains a significantly smaller number of parameters instead of all the model parameters. This produces a smaller file (~100 MBs) and makes it easier to quickly train a model to learn a new concept. LoRA weights are typically loaded into the UNet, text encoder or both. There are two classes for loading LoRA weights:
LoRA is a fast and lightweight training method that inserts and trains a significantly smaller number of parameters instead of all the model parameters. This produces a smaller file (~100 MBs) and makes it easier to quickly train a model to learn a new concept. LoRA weights are typically loaded into the denoiser, text encoder or both. The denoiser usually corresponds to a UNet ([`UNet2DConditionModel`], for example) or a Transformer ([`SD3Transformer2DModel`], for example). There are several classes for loading LoRA weights:
- [`LoraLoaderMixin`] provides functions for loading and unloading, fusing and unfusing, enabling and disabling, and more functions for managing LoRA weights. This class can be used with any model.
- [`StableDiffusionXLLoraLoaderMixin`] is a [Stable Diffusion (SDXL)](../../api/pipelines/stable_diffusion/stable_diffusion_xl) version of the [`LoraLoaderMixin`] class for loading and saving LoRA weights. It can only be used with the SDXL model.
- [`StableDiffusionLoraLoaderMixin`] provides functions for loading and unloading, fusing and unfusing, enabling and disabling, and more functions for managing LoRA weights. This class can be used with any model.
- [`StableDiffusionXLLoraLoaderMixin`] is a [Stable Diffusion (SDXL)](../../api/pipelines/stable_diffusion/stable_diffusion_xl) version of the [`StableDiffusionLoraLoaderMixin`] class for loading and saving LoRA weights. It can only be used with the SDXL model.
- [`SD3LoraLoaderMixin`] provides similar functions for [Stable Diffusion 3](https://huggingface.co/blog/sd3).
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
<Tip>
@@ -23,10 +26,22 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
</Tip>
## LoraLoaderMixin
## StableDiffusionLoraLoaderMixin
[[autodoc]] loaders.lora.LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.StableDiffusionLoraLoaderMixin
## StableDiffusionXLLoraLoaderMixin
[[autodoc]] loaders.lora.StableDiffusionXLLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.StableDiffusionXLLoraLoaderMixin
## SD3LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.SD3LoraLoaderMixin
## AmusedLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AmusedLoraLoaderMixin
## LoraBaseMixin
[[autodoc]] loaders.lora_base.LoraBaseMixin
+1 -1
View File
@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# PEFT
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`] to load an adapter.
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`], [`SD3Transformer2DModel`] to operate with an adapter.
<Tip>
@@ -22,6 +22,7 @@ The [`~loaders.FromSingleFileMixin.from_single_file`] method allows you to load:
## Supported pipelines
- [`CogVideoXPipeline`]
- [`StableDiffusionPipeline`]
- [`StableDiffusionImg2ImgPipeline`]
- [`StableDiffusionInpaintPipeline`]
@@ -49,6 +50,7 @@ The [`~loaders.FromSingleFileMixin.from_single_file`] method allows you to load:
- [`UNet2DConditionModel`]
- [`StableCascadeUNet`]
- [`AutoencoderKL`]
- [`AutoencoderKLCogVideoX`]
- [`ControlNetModel`]
- [`SD3Transformer2DModel`]
+1 -1
View File
@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# UNet
Some training methods - like LoRA and Custom Diffusion - typically target the UNet's attention layers, but these training methods can also target other non-attention layers. Instead of training all of a model's parameters, only a subset of the parameters are trained, which is faster and more efficient. This class is useful if you're *only* loading weights into a UNet. If you need to load weights into the text encoder or a text encoder and UNet, try using the [`~loaders.LoraLoaderMixin.load_lora_weights`] function instead.
Some training methods - like LoRA and Custom Diffusion - typically target the UNet's attention layers, but these training methods can also target other non-attention layers. Instead of training all of a model's parameters, only a subset of the parameters are trained, which is faster and more efficient. This class is useful if you're *only* loading weights into a UNet. If you need to load weights into the text encoder or a text encoder and UNet, try using the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] function instead.
The [`UNet2DConditionLoadersMixin`] class provides functions for loading and saving weights, fusing and unfusing LoRAs, disabling and enabling LoRAs, and setting and deleting adapters.
@@ -0,0 +1,38 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AutoencoderOobleck
The Oobleck variational autoencoder (VAE) model with KL loss was introduced in [Stability-AI/stable-audio-tools](https://github.com/Stability-AI/stable-audio-tools) and [Stable Audio Open](https://huggingface.co/papers/2407.14358) by Stability AI. The model is used in 🤗 Diffusers to encode audio waveforms into latents and to decode latent representations into audio waveforms.
The abstract from the paper is:
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
## AutoencoderOobleck
[[autodoc]] AutoencoderOobleck
- decode
- encode
- all
## OobleckDecoderOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.OobleckDecoderOutput
## OobleckDecoderOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.OobleckDecoderOutput
## AutoencoderOobleckOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.AutoencoderOobleckOutput
@@ -0,0 +1,69 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLCogVideoX
The 3D variational autoencoder (VAE) model with KL loss using CogVideoX.
## Loading from the original format
By default, the [`AutoencoderKLCogVideoX`] should be loaded with [`~ModelMixin.from_pretrained`], but it can also be loaded from the original format using [`FromOriginalModelMixin.from_single_file`] as follows:
```py
from diffusers import AutoencoderKLCogVideoX
url = "THUDM/CogVideoX-2b" # can also be a local file
model = AutoencoderKLCogVideoX.from_single_file(url)
```
## AutoencoderKLCogVideoX
[[autodoc]] AutoencoderKLCogVideoX
- decode
- encode
- all
## CogVideoXSafeConv3d
[[autodoc]] CogVideoXSafeConv3d
## CogVideoXCausalConv3d
[[autodoc]] CogVideoXCausalConv3d
## CogVideoXSpatialNorm3D
[[autodoc]] CogVideoXSpatialNorm3D
## CogVideoXResnetBlock3D
[[autodoc]] CogVideoXResnetBlock3D
## CogVideoXDownBlock3D
[[autodoc]] CogVideoXDownBlock3D
## CogVideoXMidBlock3D
[[autodoc]] CogVideoXMidBlock3D
## CogVideoXUpBlock3D
[[autodoc]] CogVideoXUpBlock3D
## CogVideoXEncoder3D
[[autodoc]] CogVideoXEncoder3D
## CogVideoXDecoder3D
[[autodoc]] CogVideoXDecoder3D
@@ -0,0 +1,18 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
## CogVideoXTransformer3DModel
A Diffusion Transformer model for 3D data from [CogVideoX](https://github.com/THUDM/CogVideoX).
## CogVideoXTransformer3DModel
[[autodoc]] CogVideoXTransformer3DModel
@@ -0,0 +1,46 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# SparseControlNetModel
SparseControlNetModel is an implementation of ControlNet for [AnimateDiff](https://arxiv.org/abs/2307.04725).
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
The SparseCtrl version of ControlNet was introduced in [SparseCtrl: Adding Sparse Controls to Text-to-Video Diffusion Models](https://arxiv.org/abs/2311.16933) for achieving controlled generation in text-to-video diffusion models by Yuwei Guo, Ceyuan Yang, Anyi Rao, Maneesh Agrawala, Dahua Lin, and Bo Dai.
The abstract from the paper is:
*The development of text-to-video (T2V), i.e., generating videos with a given text prompt, has been significantly advanced in recent years. However, relying solely on text prompts often results in ambiguous frame composition due to spatial uncertainty. The research community thus leverages the dense structure signals, e.g., per-frame depth/edge sequences, to enhance controllability, whose collection accordingly increases the burden of inference. In this work, we present SparseCtrl to enable flexible structure control with temporally sparse signals, requiring only one or a few inputs, as shown in Figure 1. It incorporates an additional condition encoder to process these sparse signals while leaving the pre-trained T2V model untouched. The proposed approach is compatible with various modalities, including sketches, depth maps, and RGB images, providing more practical control for video generation and promoting applications such as storyboarding, depth rendering, keyframe animation, and interpolation. Extensive experiments demonstrate the generalization of SparseCtrl on both original and personalized T2V generators. Codes and models will be publicly available at [this https URL](https://guoyww.github.io/projects/SparseCtrl).*
## Example for loading SparseControlNetModel
```python
import torch
from diffusers import SparseControlNetModel
# fp32 variant in float16
# 1. Scribble checkpoint
controlnet = SparseControlNetModel.from_pretrained("guoyww/animatediff-sparsectrl-scribble", torch_dtype=torch.float16)
# 2. RGB checkpoint
controlnet = SparseControlNetModel.from_pretrained("guoyww/animatediff-sparsectrl-rgb", torch_dtype=torch.float16)
# For loading fp16 variant, pass `variant="fp16"` as an additional parameter
```
## SparseControlNetModel
[[autodoc]] SparseControlNetModel
## SparseControlNetOutput
[[autodoc]] models.controlnet_sparsectrl.SparseControlNetOutput
@@ -0,0 +1,19 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# FluxTransformer2DModel
A Transformer model for image-like data from [Flux](https://blackforestlabs.ai/announcing-black-forest-labs/).
## FluxTransformer2DModel
[[autodoc]] FluxTransformer2DModel
@@ -0,0 +1,19 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# StableAudioDiTModel
A Transformer model for audio waveforms from [Stable Audio Open](https://huggingface.co/papers/2407.14358).
## StableAudioDiTModel
[[autodoc]] StableAudioDiTModel
+275 -1
View File
@@ -25,6 +25,9 @@ The abstract of the paper is the following:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
| [AnimateDiffControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_controlnet.py) | *Controlled Video-to-Video Generation with AnimateDiff using ControlNet* |
| [AnimateDiffSparseControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_sparsectrl.py) | *Controlled Video-to-Video Generation with AnimateDiff using SparseCtrl* |
| [AnimateDiffSDXLPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_sdxl.py) | *Video-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py) | *Video-to-Video Generation with AnimateDiff* |
## Available checkpoints
@@ -100,6 +103,266 @@ AnimateDiff tends to work better with finetuned Stable Diffusion models. If you
</Tip>
### AnimateDiffControlNetPipeline
AnimateDiff can also be used with ControlNets ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala. With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide depth maps, the ControlNet model generates a video that'll preserve the spatial information from the depth maps. It is a more flexible and accurate way to control the video generation process.
```python
import torch
from diffusers import AnimateDiffControlNetPipeline, AutoencoderKL, ControlNetModel, MotionAdapter, LCMScheduler
from diffusers.utils import export_to_gif, load_video
# Additionally, you will need a preprocess videos before they can be used with the ControlNet
# HF maintains just the right package for it: `pip install controlnet_aux`
from controlnet_aux.processor import ZoeDetector
# Download controlnets from https://huggingface.co/lllyasviel/ControlNet-v1-1 to use .from_single_file
# Download Diffusers-format controlnets, such as https://huggingface.co/lllyasviel/sd-controlnet-depth, to use .from_pretrained()
controlnet = ControlNetModel.from_single_file("control_v11f1p_sd15_depth.pth", torch_dtype=torch.float16)
# We use AnimateLCM for this example but one can use the original motion adapters as well (for example, https://huggingface.co/guoyww/animatediff-motion-adapter-v1-5-3)
motion_adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM")
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16)
pipe: AnimateDiffControlNetPipeline = AnimateDiffControlNetPipeline.from_pretrained(
"SG161222/Realistic_Vision_V5.1_noVAE",
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
).to(device="cuda", dtype=torch.float16)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights("wangfuyun/AnimateLCM", weight_name="AnimateLCM_sd15_t2v_lora.safetensors", adapter_name="lcm-lora")
pipe.set_adapters(["lcm-lora"], [0.8])
depth_detector = ZoeDetector.from_pretrained("lllyasviel/Annotators").to("cuda")
video = load_video("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-vid2vid-input-1.gif")
conditioning_frames = []
with pipe.progress_bar(total=len(video)) as progress_bar:
for frame in video:
conditioning_frames.append(depth_detector(frame))
progress_bar.update()
prompt = "a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality"
negative_prompt = "bad quality, worst quality"
video = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
num_frames=len(video),
num_inference_steps=10,
guidance_scale=2.0,
conditioning_frames=conditioning_frames,
generator=torch.Generator().manual_seed(42),
).frames[0]
export_to_gif(video, "animatediff_controlnet.gif", fps=8)
```
Here are some sample outputs:
<table align="center">
<tr>
<th align="center">Source Video</th>
<th align="center">Output Video</th>
</tr>
<tr>
<td align="center">
raccoon playing a guitar
<br />
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-vid2vid-input-1.gif" alt="racoon playing a guitar" />
</td>
<td align="center">
a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality
<br/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-controlnet-output.gif" alt="a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality" />
</td>
</tr>
</table>
### AnimateDiffSparseControlNetPipeline
[SparseCtrl: Adding Sparse Controls to Text-to-Video Diffusion Models](https://arxiv.org/abs/2311.16933) for achieving controlled generation in text-to-video diffusion models by Yuwei Guo, Ceyuan Yang, Anyi Rao, Maneesh Agrawala, Dahua Lin, and Bo Dai.
The abstract from the paper is:
*The development of text-to-video (T2V), i.e., generating videos with a given text prompt, has been significantly advanced in recent years. However, relying solely on text prompts often results in ambiguous frame composition due to spatial uncertainty. The research community thus leverages the dense structure signals, e.g., per-frame depth/edge sequences, to enhance controllability, whose collection accordingly increases the burden of inference. In this work, we present SparseCtrl to enable flexible structure control with temporally sparse signals, requiring only one or a few inputs, as shown in Figure 1. It incorporates an additional condition encoder to process these sparse signals while leaving the pre-trained T2V model untouched. The proposed approach is compatible with various modalities, including sketches, depth maps, and RGB images, providing more practical control for video generation and promoting applications such as storyboarding, depth rendering, keyframe animation, and interpolation. Extensive experiments demonstrate the generalization of SparseCtrl on both original and personalized T2V generators. Codes and models will be publicly available at [this https URL](https://guoyww.github.io/projects/SparseCtrl).*
SparseCtrl introduces the following checkpoints for controlled text-to-video generation:
- [SparseCtrl Scribble](https://huggingface.co/guoyww/animatediff-sparsectrl-scribble)
- [SparseCtrl RGB](https://huggingface.co/guoyww/animatediff-sparsectrl-rgb)
#### Using SparseCtrl Scribble
```python
import torch
from diffusers import AnimateDiffSparseControlNetPipeline
from diffusers.models import AutoencoderKL, MotionAdapter, SparseControlNetModel
from diffusers.schedulers import DPMSolverMultistepScheduler
from diffusers.utils import export_to_gif, load_image
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
motion_adapter_id = "guoyww/animatediff-motion-adapter-v1-5-3"
controlnet_id = "guoyww/animatediff-sparsectrl-scribble"
lora_adapter_id = "guoyww/animatediff-motion-lora-v1-5-3"
vae_id = "stabilityai/sd-vae-ft-mse"
device = "cuda"
motion_adapter = MotionAdapter.from_pretrained(motion_adapter_id, torch_dtype=torch.float16).to(device)
controlnet = SparseControlNetModel.from_pretrained(controlnet_id, torch_dtype=torch.float16).to(device)
vae = AutoencoderKL.from_pretrained(vae_id, torch_dtype=torch.float16).to(device)
scheduler = DPMSolverMultistepScheduler.from_pretrained(
model_id,
subfolder="scheduler",
beta_schedule="linear",
algorithm_type="dpmsolver++",
use_karras_sigmas=True,
)
pipe = AnimateDiffSparseControlNetPipeline.from_pretrained(
model_id,
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
scheduler=scheduler,
torch_dtype=torch.float16,
).to(device)
pipe.load_lora_weights(lora_adapter_id, adapter_name="motion_lora")
pipe.fuse_lora(lora_scale=1.0)
prompt = "an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality"
negative_prompt = "low quality, worst quality, letterboxed"
image_files = [
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-1.png",
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-2.png",
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-3.png"
]
condition_frame_indices = [0, 8, 15]
conditioning_frames = [load_image(img_file) for img_file in image_files]
video = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
num_inference_steps=25,
conditioning_frames=conditioning_frames,
controlnet_conditioning_scale=1.0,
controlnet_frame_indices=condition_frame_indices,
generator=torch.Generator().manual_seed(1337),
).frames[0]
export_to_gif(video, "output.gif")
```
Here are some sample outputs:
<table align="center">
<tr>
<center>
<b>an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality</b>
</center>
</tr>
<tr>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-1.png" alt="scribble-1" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-2.png" alt="scribble-2" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-3.png" alt="scribble-3" />
</center>
</td>
</tr>
<tr>
<td colspan=3>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-sparsectrl-scribble-results.gif" alt="an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality" />
</center>
</td>
</tr>
</table>
#### Using SparseCtrl RGB
```python
import torch
from diffusers import AnimateDiffSparseControlNetPipeline
from diffusers.models import AutoencoderKL, MotionAdapter, SparseControlNetModel
from diffusers.schedulers import DPMSolverMultistepScheduler
from diffusers.utils import export_to_gif, load_image
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
motion_adapter_id = "guoyww/animatediff-motion-adapter-v1-5-3"
controlnet_id = "guoyww/animatediff-sparsectrl-rgb"
lora_adapter_id = "guoyww/animatediff-motion-lora-v1-5-3"
vae_id = "stabilityai/sd-vae-ft-mse"
device = "cuda"
motion_adapter = MotionAdapter.from_pretrained(motion_adapter_id, torch_dtype=torch.float16).to(device)
controlnet = SparseControlNetModel.from_pretrained(controlnet_id, torch_dtype=torch.float16).to(device)
vae = AutoencoderKL.from_pretrained(vae_id, torch_dtype=torch.float16).to(device)
scheduler = DPMSolverMultistepScheduler.from_pretrained(
model_id,
subfolder="scheduler",
beta_schedule="linear",
algorithm_type="dpmsolver++",
use_karras_sigmas=True,
)
pipe = AnimateDiffSparseControlNetPipeline.from_pretrained(
model_id,
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
scheduler=scheduler,
torch_dtype=torch.float16,
).to(device)
pipe.load_lora_weights(lora_adapter_id, adapter_name="motion_lora")
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-firework.png")
video = pipe(
prompt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background",
negative_prompt="low quality, worst quality",
num_inference_steps=25,
conditioning_frames=image,
controlnet_frame_indices=[0],
controlnet_conditioning_scale=1.0,
generator=torch.Generator().manual_seed(42),
).frames[0]
export_to_gif(video, "output.gif")
```
Here are some sample outputs:
<table align="center">
<tr>
<center>
<b>closeup face photo of man in black clothes, night city street, bokeh, fireworks in background</b>
</center>
</tr>
<tr>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-firework.png" alt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-sparsectrl-rgb-result.gif" alt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background" />
</center>
</td>
</tr>
</table>
### AnimateDiffSDXLPipeline
AnimateDiff can also be used with SDXL models. This is currently an experimental feature as only a beta release of the motion adapter checkpoint is available.
@@ -571,7 +834,6 @@ ckpt_path = "https://huggingface.co/Lightricks/LongAnimateDiff/blob/main/lt_long
adapter = MotionAdapter.from_single_file(ckpt_path, torch_dtype=torch.float16)
pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapter=adapter)
```
## AnimateDiffPipeline
@@ -580,6 +842,18 @@ pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapt
- all
- __call__
## AnimateDiffControlNetPipeline
[[autodoc]] AnimateDiffControlNetPipeline
- all
- __call__
## AnimateDiffSparseControlNetPipeline
[[autodoc]] AnimateDiffSparseControlNetPipeline
- all
- __call__
## AnimateDiffSDXLPipeline
[[autodoc]] AnimateDiffSDXLPipeline
+1 -1
View File
@@ -18,7 +18,7 @@ It was developed by the Fal team and more details about it can be found in [this
<Tip>
AuraFlow can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details.
AuraFlow can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details.
</Tip>
+79
View File
@@ -0,0 +1,79 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
## TODO: The paper is still being written.
-->
# CogVideoX
[TODO]() from Tsinghua University & ZhipuAI.
The abstract from the paper is:
The paper is still being written.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
### Inference
Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) to reduce the inference latency.
First, load the pipeline:
```python
import torch
from diffusers import LattePipeline
pipeline = LattePipeline.from_pretrained(
"THUDM/CogVideoX-2b", torch_dtype=torch.float16
).to("cuda")
```
Then change the memory layout of the pipelines `transformer` and `vae` components to `torch.channels-last`:
```python
pipeline.transformer.to(memory_format=torch.channels_last)
pipeline.vae.to(memory_format=torch.channels_last)
```
Finally, compile the components and run inference:
```python
pipeline.transformer = torch.compile(pipeline.transformer)
pipeline.vae.decode = torch.compile(pipeline.vae.decode)
# CogVideoX works very well with long and well-described prompts
prompt = "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical atmosphere of this unique musical performance."
video = pipeline(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
```
The [benchmark](TODO: link) results on an 80GB A100 machine are:
```
Without torch.compile(): Average inference time: TODO seconds.
With torch.compile(): Average inference time: TODO seconds.
```
## CogVideoXPipeline
[[autodoc]] CogVideoXPipeline
- all
- __call__
## CogVideoXPipelineOutput
[[autodoc]] pipelines.pipline_cogvideo.pipeline_output.CogVideoXPipelineOutput
+84
View File
@@ -0,0 +1,84 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Flux
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
<Tip>
Flux can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details. Additionally, Flux can benefit from quantization for memory efficiency with a trade-off in inference latency. Refer to [this blog post](https://huggingface.co/blog/quanto-diffusers) to learn more. For an exhaustive list of resources, check out [this gist](https://gist.github.com/sayakpaul/b664605caf0aa3bf8585ab109dd5ac9c).
</Tip>
Flux comes in two variants:
* Timestep-distilled (`black-forest-labs/FLUX.1-schnell`)
* Guidance-distilled (`black-forest-labs/FLUX.1-dev`)
Both checkpoints have slightly difference usage which we detail below.
### Timestep-distilled
* `max_sequence_length` cannot be more than 256.
* `guidance_scale` needs to be 0.
* As this is a timestep-distilled model, it benefits from fewer sampling steps.
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-schnell", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "A cat holding a sign that says hello world"
out = pipe(
prompt=prompt,
guidance_scale=0.,
height=768,
width=1360,
num_inference_steps=4,
max_sequence_length=256,
).images[0]
out.save("image.png")
```
### Guidance-distilled
* The guidance-distilled variant takes about 50 sampling steps for good-quality generation.
* It doesn't have any limitations around the `max_sequence_length`.
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "a tiny astronaut hatching from an egg on the moon"
out = pipe(
prompt=prompt,
guidance_scale=3.5,
height=768,
width=1360,
num_inference_steps=50,
).images[0]
out.save("image.png")
```
## FluxPipeline
[[autodoc]] FluxPipeline
- all
- __call__
+58
View File
@@ -41,6 +41,64 @@ image = pipe(
image.save("kolors_sample.png")
```
### IP Adapter
Kolors needs a different IP Adapter to work, and it uses [Openai-CLIP-336](https://huggingface.co/openai/clip-vit-large-patch14-336) as an image encoder.
<Tip>
Using an IP Adapter with Kolors requires more than 24GB of VRAM. To use it, we recommend using [`~DiffusionPipeline.enable_model_cpu_offload`] on consumer GPUs.
</Tip>
<Tip>
While Kolors is integrated in Diffusers, you need to load the image encoder from a revision to use the safetensor files. You can still use the main branch of the original repository if you're comfortable loading pickle checkpoints.
</Tip>
```python
import torch
from transformers import CLIPVisionModelWithProjection
from diffusers import DPMSolverMultistepScheduler, KolorsPipeline
from diffusers.utils import load_image
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"Kwai-Kolors/Kolors-IP-Adapter-Plus",
subfolder="image_encoder",
low_cpu_mem_usage=True,
torch_dtype=torch.float16,
revision="refs/pr/4",
)
pipe = KolorsPipeline.from_pretrained(
"Kwai-Kolors/Kolors-diffusers", image_encoder=image_encoder, torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config, use_karras_sigmas=True)
pipe.load_ip_adapter(
"Kwai-Kolors/Kolors-IP-Adapter-Plus",
subfolder="",
weight_name="ip_adapter_plus_general.safetensors",
revision="refs/pr/4",
image_encoder_folder=None,
)
pipe.enable_model_cpu_offload()
ipa_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/cat_square.png")
image = pipe(
prompt="best quality, high quality",
negative_prompt="",
guidance_scale=6.5,
num_inference_steps=25,
ip_adapter_image=ipa_image,
).images[0]
image.save("kolors_ipa_sample.png")
```
## KolorsPipeline
[[autodoc]] KolorsPipeline
+2
View File
@@ -24,6 +24,8 @@ The abstract from the paper is:
**Highlights**: Latte is a latent diffusion transformer proposed as a backbone for modeling different modalities (trained for text-to-video generation here). It achieves state-of-the-art performance across four standard video benchmarks - [FaceForensics](https://arxiv.org/abs/1803.09179), [SkyTimelapse](https://arxiv.org/abs/1709.07592), [UCF101](https://arxiv.org/abs/1212.0402) and [Taichi-HD](https://arxiv.org/abs/2003.00196). To prepare and download the datasets for evaluation, please refer to [this https URL](https://github.com/Vchitect/Latte/blob/main/docs/datasets_evaluation.md).
This pipeline was contributed by [maxin-cn](https://github.com/maxin-cn). The original codebase can be found [here](https://github.com/Vchitect/Latte). The original weights can be found under [hf.co/maxin-cn](https://huggingface.co/maxin-cn).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
+4 -2
View File
@@ -43,6 +43,8 @@ Lumina-T2X has the following components:
* It uses a Flow-based Large Diffusion Transformer as the backbone
* It supports different any modalities with one backbone and corresponding encoder, decoder.
This pipeline was contributed by [PommesPeter](https://github.com/PommesPeter). The original codebase can be found [here](https://github.com/Alpha-VLLM/Lumina-T2X). The original weights can be found under [hf.co/Alpha-VLLM](https://huggingface.co/Alpha-VLLM).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
@@ -57,7 +59,7 @@ First, load the pipeline:
```python
from diffusers import LuminaText2ImgPipeline
import torch
import torch
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers", torch_dtype=torch.bfloat16
@@ -85,4 +87,4 @@ image = pipeline(prompt="Upper body of a young woman in a Victorian-era outfit w
[[autodoc]] LuminaText2ImgPipeline
- all
- __call__
+1
View File
@@ -71,6 +71,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Semantic Guidance](semantic_stable_diffusion) | text2image |
| [Shap-E](shap_e) | text-to-3D, image-to-3D |
| [Spectrogram Diffusion](spectrogram_diffusion) | |
| [Stable Audio](stable_audio) | text2audio |
| [Stable Diffusion](stable_diffusion/overview) | text2image, image2image, depth2image, inpainting, image variation, latent upscaler, super-resolution |
| [Stable Diffusion Model Editing](model_editing) | model editing |
| [Stable Diffusion XL](stable_diffusion/stable_diffusion_xl) | text2image, image2image, inpainting |
+29
View File
@@ -20,6 +20,29 @@ The abstract from the paper is:
*Recent studies have demonstrated that diffusion models are capable of generating high-quality samples, but their quality heavily depends on sampling guidance techniques, such as classifier guidance (CG) and classifier-free guidance (CFG). These techniques are often not applicable in unconditional generation or in various downstream tasks such as image restoration. In this paper, we propose a novel sampling guidance, called Perturbed-Attention Guidance (PAG), which improves diffusion sample quality across both unconditional and conditional settings, achieving this without requiring additional training or the integration of external modules. PAG is designed to progressively enhance the structure of samples throughout the denoising process. It involves generating intermediate samples with degraded structure by substituting selected self-attention maps in diffusion U-Net with an identity matrix, by considering the self-attention mechanisms' ability to capture structural information, and guiding the denoising process away from these degraded samples. In both ADM and Stable Diffusion, PAG surprisingly improves sample quality in conditional and even unconditional scenarios. Moreover, PAG significantly improves the baseline performance in various downstream tasks where existing guidances such as CG or CFG cannot be fully utilized, including ControlNet with empty prompts and image restoration such as inpainting and deblurring.*
PAG can be used by specifying the `pag_applied_layers` as a parameter when instantiating a PAG pipeline. It can be a single string or a list of strings. Each string can be a unique layer identifier or a regular expression to identify one or more layers.
- Full identifier as a normal string: `down_blocks.2.attentions.0.transformer_blocks.0.attn1.processor`
- Full identifier as a RegEx: `down_blocks.2.(attentions|motion_modules).0.transformer_blocks.0.attn1.processor`
- Partial identifier as a RegEx: `down_blocks.2`, or `attn1`
- List of identifiers (can be combo of strings and ReGex): `["blocks.1", "blocks.(14|20)", r"down_blocks\.(2,3)"]`
<Tip warning={true}>
Since RegEx is supported as a way for matching layer identifiers, it is crucial to use it correctly otherwise there might be unexpected behaviour. The recommended way to use PAG is by specifying layers as `blocks.{layer_index}` and `blocks.({layer_index_1|layer_index_2|...})`. Using it in any other way, while doable, may bypass our basic validation checks and give you unexpected results.
</Tip>
## AnimateDiffPAGPipeline
[[autodoc]] AnimateDiffPAGPipeline
- all
- __call__
## HunyuanDiTPAGPipeline
[[autodoc]] HunyuanDiTPAGPipeline
- all
- __call__
## StableDiffusionPAGPipeline
[[autodoc]] StableDiffusionPAGPipeline
- all
@@ -49,3 +72,9 @@ The abstract from the paper is:
[[autodoc]] StableDiffusionXLControlNetPAGPipeline
- all
- __call__
## PixArtSigmaPAGPipeline
[[autodoc]] PixArtSigmaPAGPipeline
- all
- __call__
@@ -0,0 +1,42 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Audio
Stable Audio was proposed in [Stable Audio Open](https://arxiv.org/abs/2407.14358) by Zach Evans et al. . it takes a text prompt as input and predicts the corresponding sound or music sample.
Stable Audio Open generates variable-length (up to 47s) stereo audio at 44.1kHz from text prompts. It comprises three components: an autoencoder that compresses waveforms into a manageable sequence length, a T5-based text embedding for text conditioning, and a transformer-based diffusion (DiT) model that operates in the latent space of the autoencoder.
Stable Audio is trained on a corpus of around 48k audio recordings, where around 47k are from Freesound and the rest are from the Free Music Archive (FMA). All audio files are licensed under CC0, CC BY, or CC Sampling+. This data is used to train the autoencoder and the DiT.
The abstract of the paper is the following:
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tool](https://github.com/Stability-AI/stable-audio-tool).
## Tips
When constructing a prompt, keep in mind:
* Descriptive prompt inputs work best; use adjectives to describe the sound (for example, "high quality" or "clear") and make the prompt context specific where possible (e.g. "melodic techno with a fast beat and synths" works better than "techno").
* Using a *negative prompt* can significantly improve the quality of the generated audio. Try using a negative prompt of "low quality, average quality".
During inference:
* The _quality_ of the generated audio sample can be controlled by the `num_inference_steps` argument; higher steps give higher quality audio at the expense of slower inference.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1 to enable. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
## StableAudioPipeline
[[autodoc]] StableAudioPipeline
- all
- __call__
@@ -0,0 +1,24 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CosineDPMSolverMultistepScheduler
The [`CosineDPMSolverMultistepScheduler`] is a variant of [`DPMSolverMultistepScheduler`] with cosine schedule, proposed by Nichol and Dhariwal (2021).
It is being used in the [Stable Audio Open](https://arxiv.org/abs/2407.14358) paper and the [Stability-AI/stable-audio-tool](https://github.com/Stability-AI/stable-audio-tool) codebase.
This scheduler was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe).
## CosineDPMSolverMultistepScheduler
[[autodoc]] CosineDPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
+2 -2
View File
@@ -340,8 +340,8 @@ Now you can wrap all these components together in a training loop with 🤗 Acce
... loss = F.mse_loss(noise_pred, noise)
... accelerator.backward(loss)
... if (step + 1) % config.gradient_accumulation_steps == 0:
... accelerator.clip_grad_norm_(model.parameters(), 1.0)
... if accelerator.sync_gradients:
... accelerator.clip_grad_norm_(model.parameters(), 1.0)
... optimizer.step()
... lr_scheduler.step()
... optimizer.zero_grad()
+2 -2
View File
@@ -35,7 +35,7 @@ pip3 install --pre torch --index-url https://download.pytorch.org/whl/nightly/cu
```
> [!TIP]
> The results reported below are from a 80GB 400W A100 with its clock rate set to the maximum.
> The results reported below are from a 80GB 400W A100 with its clock rate set to the maximum.
> If you're interested in the full benchmarking code, take a look at [huggingface/diffusion-fast](https://github.com/huggingface/diffusion-fast).
@@ -168,7 +168,7 @@ Using SDPA attention and compiling both the UNet and VAE cuts the latency from 3
</div>
> [!TIP]
> From PyTorch 2.3.1, you can control the caching behavior of `torch.compile()`. This is particularly beneficial for compilation modes like `"max-autotune"` which performs a grid-search over several compilation flags to find the optimal configuration. Learn more in the [Compile Time Caching in torch.compile](https://pytorch.org/tutorials/recipes/torch_compile_caching_tutorial.html) tutorial.
> From PyTorch 2.3.1, you can control the caching behavior of `torch.compile()`. This is particularly beneficial for compilation modes like `"max-autotune"` which performs a grid-search over several compilation flags to find the optimal configuration. Learn more in the [Compile Time Caching in torch.compile](https://pytorch.org/tutorials/recipes/torch_compile_caching_tutorial.html) tutorial.
### Prevent graph breaks
@@ -18,13 +18,13 @@ A modern diffusion model, like [Stable Diffusion XL (SDXL)](../using-diffusers/s
* Two text encoders
* A UNet for denoising
Usually, the text encoders and the denoiser are much larger compared to the VAE.
Usually, the text encoders and the denoiser are much larger compared to the VAE.
As models get bigger and better, its possible your model is so big that even a single copy wont fit in memory. But that doesnt mean it cant be loaded. If you have more than one GPU, there is more memory available to store your model. In this case, its better to split your model checkpoint into several smaller *checkpoint shards*.
When a text encoder checkpoint has multiple shards, like [T5-xxl for SD3](https://huggingface.co/stabilityai/stable-diffusion-3-medium-diffusers/tree/main/text_encoder_3), it is automatically handled by the [Transformers](https://huggingface.co/docs/transformers/index) library as it is a required dependency of Diffusers when using the [`StableDiffusion3Pipeline`]. More specifically, Transformers will automatically handle the loading of multiple shards within the requested model class and get it ready so that inference can be performed.
The denoiser checkpoint can also have multiple shards and supports inference thanks to the [Accelerate](https://huggingface.co/docs/accelerate/index) library.
The denoiser checkpoint can also have multiple shards and supports inference thanks to the [Accelerate](https://huggingface.co/docs/accelerate/index) library.
> [!TIP]
> Refer to the [Handling big models for inference](https://huggingface.co/docs/accelerate/main/en/concept_guides/big_model_inference) guide for general guidance when working with big models that are hard to fit into memory.
@@ -43,7 +43,7 @@ unet.save_pretrained("sdxl-unet-sharded", max_shard_size="5GB")
The size of the fp32 variant of the SDXL UNet checkpoint is ~10.4GB. Set the `max_shard_size` parameter to 5GB to create 3 shards. After saving, you can load them in [`StableDiffusionXLPipeline`]:
```python
from diffusers import UNet2DConditionModel, StableDiffusionXLPipeline
from diffusers import UNet2DConditionModel, StableDiffusionXLPipeline
import torch
unet = UNet2DConditionModel.from_pretrained(
@@ -57,14 +57,14 @@ image = pipeline("a cute dog running on the grass", num_inference_steps=30).imag
image.save("dog.png")
```
If placing all the model-level components on the GPU at once is not feasible, use [`~DiffusionPipeline.enable_model_cpu_offload`] to help you:
If placing all the model-level components on the GPU at once is not feasible, use [`~DiffusionPipeline.enable_model_cpu_offload`] to help you:
```diff
- pipeline.to("cuda")
+ pipeline.enable_model_cpu_offload()
```
In general, we recommend sharding when a checkpoint is more than 5GB (in fp32).
In general, we recommend sharding when a checkpoint is more than 5GB (in fp32).
## Device placement
@@ -191,7 +191,7 @@ image
## Manage active adapters
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, use the [`~diffusers.loaders.LoraLoaderMixin.get_active_adapters`] method to check the list of active adapters:
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, use the [`~diffusers.loaders.StableDiffusionLoraLoaderMixin.get_active_adapters`] method to check the list of active adapters:
```py
active_adapters = pipe.get_active_adapters()
@@ -199,7 +199,7 @@ active_adapters
["toy", "pixel"]
```
You can also get the active adapters of each pipeline component with [`~diffusers.loaders.LoraLoaderMixin.get_list_adapters`]:
You can also get the active adapters of each pipeline component with [`~diffusers.loaders.StableDiffusionLoraLoaderMixin.get_list_adapters`]:
```py
list_adapters_component_wise = pipe.get_list_adapters()
+1 -1
View File
@@ -256,7 +256,7 @@ make_image_grid([init_image, mask_image, output], rows=1, cols=3)
## Guess mode
[Guess mode](https://github.com/lllyasviel/ControlNet/discussions/188) does not require supplying a prompt to a ControlNet at all! This forces the ControlNet encoder to do it's best to "guess" the contents of the input control map (depth map, pose estimation, canny edge, etc.).
[Guess mode](https://github.com/lllyasviel/ControlNet/discussions/188) does not require supplying a prompt to a ControlNet at all! This forces the ControlNet encoder to do its best to "guess" the contents of the input control map (depth map, pose estimation, canny edge, etc.).
Guess mode adjusts the scale of the output residuals from a ControlNet by a fixed ratio depending on the block depth. The shallowest `DownBlock` corresponds to 0.1, and as the blocks get deeper, the scale increases exponentially such that the scale of the `MidBlock` output becomes 1.0.
@@ -64,7 +64,7 @@ image
</hfoption>
<hfoption id="LCM-LoRA">
To use LCM-LoRAs, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt to generate an image in just 4 steps.
To use LCM-LoRAs, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt to generate an image in just 4 steps.
A couple of notes to keep in mind when using LCM-LoRAs are:
@@ -156,7 +156,7 @@ image
</hfoption>
<hfoption id="LCM-LoRA">
To use LCM-LoRAs for image-to-image, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt and initial image to generate an image in just 4 steps.
To use LCM-LoRAs for image-to-image, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt and initial image to generate an image in just 4 steps.
> [!TIP]
> Experiment with different values for `num_inference_steps`, `strength`, and `guidance_scale` to get the best results.
@@ -207,7 +207,7 @@ image
## Inpainting
To use LCM-LoRAs for inpainting, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt, initial image, and mask image to generate an image in just 4 steps.
To use LCM-LoRAs for inpainting, you need to replace the scheduler with the [`LCMScheduler`] and load the LCM-LoRA weights with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. Then you can use the pipeline as usual, and pass a text prompt, initial image, and mask image to generate an image in just 4 steps.
```py
import torch
@@ -262,7 +262,7 @@ LCMs are compatible with adapters like LoRA, ControlNet, T2I-Adapter, and Animat
<hfoptions id="lcm-lora">
<hfoption id="LCM">
Load the LCM checkpoint for your supported model into [`UNet2DConditionModel`] and replace the scheduler with the [`LCMScheduler`]. Then you can use the [`~loaders.LoraLoaderMixin.load_lora_weights`] method to load the LoRA weights into the LCM and generate a styled image in a few steps.
Load the LCM checkpoint for your supported model into [`UNet2DConditionModel`] and replace the scheduler with the [`LCMScheduler`]. Then you can use the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method to load the LoRA weights into the LCM and generate a styled image in a few steps.
```python
from diffusers import StableDiffusionXLPipeline, UNet2DConditionModel, LCMScheduler
@@ -294,7 +294,7 @@ image
</hfoption>
<hfoption id="LCM-LoRA">
Replace the scheduler with the [`LCMScheduler`]. Then you can use the [`~loaders.LoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights and the style LoRA you want to use. Combine both LoRA adapters with the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method and generate a styled image in a few steps.
Replace the scheduler with the [`LCMScheduler`]. Then you can use the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights and the style LoRA you want to use. Combine both LoRA adapters with the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method and generate a styled image in a few steps.
```py
import torch
@@ -389,7 +389,7 @@ make_image_grid([canny_image, image], rows=1, cols=2)
</hfoption>
<hfoption id="LCM-LoRA">
Load a ControlNet model trained on canny images and pass it to the [`ControlNetModel`]. Then you can load a Stable Diffusion v1.5 model into [`StableDiffusionControlNetPipeline`] and replace the scheduler with the [`LCMScheduler`]. Use the [`~loaders.LoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights, and pass the canny image to the pipeline and generate an image.
Load a ControlNet model trained on canny images and pass it to the [`ControlNetModel`]. Then you can load a Stable Diffusion v1.5 model into [`StableDiffusionControlNetPipeline`] and replace the scheduler with the [`LCMScheduler`]. Use the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights, and pass the canny image to the pipeline and generate an image.
> [!TIP]
> Experiment with different values for `num_inference_steps`, `controlnet_conditioning_scale`, `cross_attention_kwargs`, and `guidance_scale` to get the best results.
@@ -525,7 +525,7 @@ image = pipe(
</hfoption>
<hfoption id="LCM-LoRA">
Load a T2IAdapter trained on canny images and pass it to the [`StableDiffusionXLAdapterPipeline`]. Replace the scheduler with the [`LCMScheduler`], and use the [`~loaders.LoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights. Pass the canny image to the pipeline and generate an image.
Load a T2IAdapter trained on canny images and pass it to the [`StableDiffusionXLAdapterPipeline`]. Replace the scheduler with the [`LCMScheduler`], and use the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method to load the LCM-LoRA weights. Pass the canny image to the pipeline and generate an image.
```py
import torch
@@ -116,7 +116,7 @@ import torch
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
```
Then use the [`~loaders.LoraLoaderMixin.load_lora_weights`] method to load the [ostris/super-cereal-sdxl-lora](https://huggingface.co/ostris/super-cereal-sdxl-lora) weights and specify the weights filename from the repository:
Then use the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method to load the [ostris/super-cereal-sdxl-lora](https://huggingface.co/ostris/super-cereal-sdxl-lora) weights and specify the weights filename from the repository:
```py
pipeline.load_lora_weights("ostris/super-cereal-sdxl-lora", weight_name="cereal_box_sdxl_v1.safetensors")
@@ -129,7 +129,7 @@ image
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_lora.png" />
</div>
The [`~loaders.LoraLoaderMixin.load_lora_weights`] method loads LoRA weights into both the UNet and text encoder. It is the preferred way for loading LoRAs because it can handle cases where:
The [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method loads LoRA weights into both the UNet and text encoder. It is the preferred way for loading LoRAs because it can handle cases where:
- the LoRA weights don't have separate identifiers for the UNet and text encoder
- the LoRA weights have separate identifiers for the UNet and text encoder
@@ -153,7 +153,7 @@ image
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_attn_proc.png" />
</div>
To unload the LoRA weights, use the [`~loaders.LoraLoaderMixin.unload_lora_weights`] method to discard the LoRA weights and restore the model to its original weights:
To unload the LoRA weights, use the [`~loaders.StableDiffusionLoraLoaderMixin.unload_lora_weights`] method to discard the LoRA weights and restore the model to its original weights:
```py
pipeline.unload_lora_weights()
@@ -161,9 +161,9 @@ pipeline.unload_lora_weights()
### Adjust LoRA weight scale
For both [`~loaders.LoraLoaderMixin.load_lora_weights`] and [`~loaders.UNet2DConditionLoadersMixin.load_attn_procs`], you can pass the `cross_attention_kwargs={"scale": 0.5}` parameter to adjust how much of the LoRA weights to use. A value of `0` is the same as only using the base model weights, and a value of `1` is equivalent to using the fully finetuned LoRA.
For both [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] and [`~loaders.UNet2DConditionLoadersMixin.load_attn_procs`], you can pass the `cross_attention_kwargs={"scale": 0.5}` parameter to adjust how much of the LoRA weights to use. A value of `0` is the same as only using the base model weights, and a value of `1` is equivalent to using the fully finetuned LoRA.
For more granular control on the amount of LoRA weights used per layer, you can use [`~loaders.LoraLoaderMixin.set_adapters`] and pass a dictionary specifying by how much to scale the weights in each layer by.
For more granular control on the amount of LoRA weights used per layer, you can use [`~loaders.StableDiffusionLoraLoaderMixin.set_adapters`] and pass a dictionary specifying by how much to scale the weights in each layer by.
```python
pipe = ... # create pipeline
pipe.load_lora_weights(..., adapter_name="my_adapter")
@@ -186,7 +186,7 @@ This also works with multiple adapters - see [this guide](https://huggingface.co
<Tip warning={true}>
Currently, [`~loaders.LoraLoaderMixin.set_adapters`] only supports scaling attention weights. If a LoRA has other parts (e.g., resnets or down-/upsamplers), they will keep a scale of 1.0.
Currently, [`~loaders.StableDiffusionLoraLoaderMixin.set_adapters`] only supports scaling attention weights. If a LoRA has other parts (e.g., resnets or down-/upsamplers), they will keep a scale of 1.0.
</Tip>
@@ -203,7 +203,7 @@ To load a Kohya LoRA, let's download the [Blueprintify SD XL 1.0](https://civita
!wget https://civitai.com/api/download/models/168776 -O blueprintify-sd-xl-10.safetensors
```
Load the LoRA checkpoint with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method, and specify the filename in the `weight_name` parameter:
Load the LoRA checkpoint with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method, and specify the filename in the `weight_name` parameter:
```py
from diffusers import AutoPipelineForText2Image
@@ -227,7 +227,7 @@ image
Some limitations of using Kohya LoRAs with 🤗 Diffusers include:
- Images may not look like those generated by UIs - like ComfyUI - for multiple reasons, which are explained [here](https://github.com/huggingface/diffusers/pull/4287/#issuecomment-1655110736).
- [LyCORIS checkpoints](https://github.com/KohakuBlueleaf/LyCORIS) aren't fully supported. The [`~loaders.LoraLoaderMixin.load_lora_weights`] method loads LyCORIS checkpoints with LoRA and LoCon modules, but Hada and LoKR are not supported.
- [LyCORIS checkpoints](https://github.com/KohakuBlueleaf/LyCORIS) aren't fully supported. The [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method loads LyCORIS checkpoints with LoRA and LoCon modules, but Hada and LoKR are not supported.
</Tip>
@@ -14,9 +14,9 @@ specific language governing permissions and limitations under the License.
It can be fun and creative to use multiple [LoRAs]((https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora)) together to generate something entirely new and unique. This works by merging multiple LoRA weights together to produce images that are a blend of different styles. Diffusers provides a few methods to merge LoRAs depending on *how* you want to merge their weights, which can affect image quality.
This guide will show you how to merge LoRAs using the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.LoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
This guide will show you how to merge LoRAs using the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style]() and [Norod78/sdxl-chalkboarddrawing-lora]() LoRAs with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style]() and [Norod78/sdxl-chalkboarddrawing-lora]() LoRAs with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
```py
from diffusers import DiffusionPipeline
@@ -182,9 +182,9 @@ image
## fuse_lora
Both the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.LoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
Both the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.LoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
For example, if you have a base model and adapters loaded and set as active with the following adapter weights:
@@ -199,13 +199,13 @@ pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
```
Fuse these LoRAs into the UNet with the [`~loaders.LoraLoaderMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.LoraLoaderMixin.fuse_lora`] method because it wont work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
Fuse these LoRAs into the UNet with the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method because it wont work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
```py
pipeline.fuse_lora(adapter_names=["ikea", "feng"], lora_scale=1.0)
```
Then you should use [`~loaders.LoraLoaderMixin.unload_lora_weights`] to unload the LoRA weights since they've already been fused with the underlying base model. Finally, call [`~DiffusionPipeline.save_pretrained`] to save the fused pipeline locally or you could call [`~DiffusionPipeline.push_to_hub`] to push the fused pipeline to the Hub.
Then you should use [`~loaders.StableDiffusionLoraLoaderMixin.unload_lora_weights`] to unload the LoRA weights since they've already been fused with the underlying base model. Finally, call [`~DiffusionPipeline.save_pretrained`] to save the fused pipeline locally or you could call [`~DiffusionPipeline.push_to_hub`] to push the fused pipeline to the Hub.
```py
pipeline.unload_lora_weights()
@@ -226,7 +226,7 @@ image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai"
image
```
You can call [`~loaders.LoraLoaderMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
You can call [`~loaders.StableDiffusionLoraLoaderMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
```py
pipeline.unfuse_lora()
@@ -74,7 +74,7 @@ pipeline = StableDiffusionPipeline.from_single_file(
[LoRA](https://hf.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) is a lightweight adapter that is fast and easy to train, making them especially popular for generating images in a certain way or style. These adapters are commonly stored in a safetensors file, and are widely popular on model sharing platforms like [civitai](https://civitai.com/).
LoRAs are loaded into a base model with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method.
LoRAs are loaded into a base model with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method.
```py
from diffusers import StableDiffusionXLPipeline
+9 -9
View File
@@ -22,7 +22,7 @@ This guide will show you how to use PAG for various tasks and use cases.
You can apply PAG to the [`StableDiffusionXLPipeline`] for tasks such as text-to-image, image-to-image, and inpainting. To enable PAG for a specific task, load the pipeline using the [AutoPipeline](../api/pipelines/auto_pipeline) API with the `enable_pag=True` flag and the `pag_applied_layers` argument.
> [!TIP]
> 🤗 Diffusers currently only supports using PAG with selected SDXL pipelines, but feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you want to add PAG support to a new pipeline!
> 🤗 Diffusers currently only supports using PAG with selected SDXL pipelines and [`PixArtSigmaPAGPipeline`]. But feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you want to add PAG support to a new pipeline!
<hfoptions id="tasks">
<hfoption id="Text-to-image">
@@ -130,10 +130,10 @@ prompt = "a dog catching a frisbee in the jungle"
generator = torch.Generator(device="cpu").manual_seed(0)
image = pipeline(
prompt,
image=init_image,
strength=0.8,
guidance_scale=guidance_scale,
prompt,
image=init_image,
strength=0.8,
guidance_scale=guidance_scale,
pag_scale=pag_scale,
generator=generator).images[0]
```
@@ -161,14 +161,14 @@ pipeline_inpaint = AutoPipelineForInpaiting.from_pretrained("stabilityai/stable-
pipeline = AutoPipelineForInpaiting.from_pipe(pipeline_inpaint, enable_pag=True)
```
This still works when your pipeline has a different task:
This still works when your pipeline has a different task:
```py
pipeline_t2i = AutoPipelineForText2Image.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16)
pipeline = AutoPipelineForInpaiting.from_pipe(pipeline_t2i, enable_pag=True)
```
Let's generate an image!
Let's generate an image!
```py
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
@@ -258,7 +258,7 @@ for pag_scale in [0.0, 3.0]:
</div>
</div>
## PAG with IP-Adapter
## PAG with IP-Adapter
[IP-Adapter](https://hf.co/papers/2308.06721) is a popular model that can be plugged into diffusion models to enable image prompting without any changes to the underlying model. You can enable PAG on a pipeline with IP-Adapter loaded.
@@ -317,7 +317,7 @@ PAG reduces artifacts and improves the overall compposition.
</div>
## Configure parameters
## Configure parameters
### pag_applied_layers
+1 -1
View File
@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# 철학 [[philosophy]]
# 철학 [[philosophy]]
🧨 Diffusers는 다양한 모달리티에서 **최신의** 사전 훈련된 diffusion 모델을 제공합니다.
그 목적은 추론과 훈련을 위한 **모듈식 툴박스**로 사용되는 것입니다.
@@ -127,7 +127,7 @@ image = pipeline(prompt, num_inference_steps=50).images[0]
[Automatic1111](https://github.com/AUTOMATIC1111/stable-diffusion-webui) (A1111)은 Stable Diffusion을 위해 널리 사용되는 웹 UI로, [Civitai](https://civitai.com/) 와 같은 모델 공유 플랫폼을 지원합니다. 특히 LoRA 기법으로 학습된 모델은 학습 속도가 빠르고 완전히 파인튜닝된 모델보다 파일 크기가 훨씬 작기 때문에 인기가 높습니다.
🤗 Diffusers는 [`~loaders.LoraLoaderMixin.load_lora_weights`]:를 사용하여 A1111 LoRA 체크포인트 불러오기를 지원합니다:
🤗 Diffusers는 [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`]:를 사용하여 A1111 LoRA 체크포인트 불러오기를 지원합니다:
```py
from diffusers import DiffusionPipeline, UniPCMultistepScheduler
+2 -2
View File
@@ -52,7 +52,7 @@ pipeline = pipeline.to("cuda")
Text-to-image의 경우 텍스트 프롬프트를 전달합니다. 기본적으로 SDXL Turbo는 512x512 이미지를 생성하며, 이 해상도에서 최상의 결과를 제공합니다. `height``width` 매개 변수를 768x768 또는 1024x1024로 설정할 수 있지만 이 경우 품질 저하를 예상할 수 있습니다.
모델이 `guidance_scale` 없이 학습되었으므로 이를 0.0으로 설정해 비활성화해야 합니다. 단일 추론 스텝만으로도 고품질 이미지를 생성할 수 있습니다.
모델이 `guidance_scale` 없이 학습되었으므로 이를 0.0으로 설정해 비활성화해야 합니다. 단일 추론 스텝만으로도 고품질 이미지를 생성할 수 있습니다.
스텝 수를 2, 3 또는 4로 늘리면 이미지 품질이 향상됩니다.
```py
@@ -74,7 +74,7 @@ image
## Image-to-image
Image-to-image 생성의 경우 `num_inference_steps * strength`가 1보다 크거나 같은지 확인하세요.
Image-to-image 생성의 경우 `num_inference_steps * strength`가 1보다 크거나 같은지 확인하세요.
Image-to-image 파이프라인은 아래 예제에서 `0.5 * 2.0 = 1` 스텝과 같이 `int(num_inference_steps * strength)` 스텝으로 실행됩니다.
```py
+1 -1
View File
@@ -21,7 +21,7 @@ specific language governing permissions and limitations under the License.
시작하기 전에 다음 라이브러리가 설치되어 있는지 확인하세요:
```py
!pip install -q -U diffusers transformers accelerate
!pip install -q -U diffusers transformers accelerate
```
이 모델에는 [SVD](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid)와 [SVD-XT](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid-xt) 두 가지 종류가 있습니다. SVD 체크포인트는 14개의 프레임을 생성하도록 학습되었고, SVD-XT 체크포인트는 25개의 프레임을 생성하도록 파인튜닝되었습니다.
@@ -57,7 +57,7 @@ from diffusers import (
StableDiffusionPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import compute_snr
from diffusers.utils import (
@@ -1318,11 +1318,11 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
LoraLoaderMixin.load_lora_into_unet(lora_state_dict, network_alphas=network_alphas, unet=unet_)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
StableDiffusionLoraLoaderMixin.load_lora_into_unet(lora_state_dict, network_alphas=network_alphas, unet=unet_)
text_encoder_state_dict = {k: v for k, v in lora_state_dict.items() if "text_encoder." in k}
LoraLoaderMixin.load_lora_into_text_encoder(
StableDiffusionLoraLoaderMixin.load_lora_into_text_encoder(
text_encoder_state_dict, network_alphas=network_alphas, text_encoder=text_encoder_one_
)
@@ -60,7 +60,7 @@ from diffusers import (
StableDiffusionXLPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params, compute_snr
from diffusers.utils import (
@@ -1646,7 +1646,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
+5 -5
View File
@@ -41,7 +41,7 @@ from transformers import (
import diffusers.optimization
from diffusers import AmusedPipeline, AmusedScheduler, EMAModel, UVit2DModel, VQModel
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import AmusedLoraLoaderMixin
from diffusers.utils import is_wandb_available
@@ -532,7 +532,7 @@ def main(args):
weights.pop()
if transformer_lora_layers_to_save is not None or text_encoder_lora_layers_to_save is not None:
LoraLoaderMixin.save_lora_weights(
AmusedLoraLoaderMixin.save_lora_weights(
output_dir,
transformer_lora_layers=transformer_lora_layers_to_save,
text_encoder_lora_layers=text_encoder_lora_layers_to_save,
@@ -566,11 +566,11 @@ def main(args):
raise ValueError(f"unexpected save model: {model.__class__}")
if transformer is not None or text_encoder_ is not None:
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
LoraLoaderMixin.load_lora_into_text_encoder(
lora_state_dict, network_alphas = AmusedLoraLoaderMixin.lora_state_dict(input_dir)
AmusedLoraLoaderMixin.load_lora_into_text_encoder(
lora_state_dict, network_alphas=network_alphas, text_encoder=text_encoder_
)
LoraLoaderMixin.load_lora_into_transformer(
AmusedLoraLoaderMixin.load_lora_into_transformer(
lora_state_dict, network_alphas=network_alphas, transformer=transformer
)
+15 -16
View File
@@ -1487,17 +1487,16 @@ NOTE: The ONNX conversions and TensorRT engine build may take up to 30 minutes.
```python
import torch
from diffusers import DDIMScheduler
from diffusers.pipelines.stable_diffusion import StableDiffusionPipeline
from diffusers.pipelines import DiffusionPipeline
# Use the DDIMScheduler scheduler here instead
scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1",
subfolder="scheduler")
scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1", subfolder="scheduler")
pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
custom_pipeline="stable_diffusion_tensorrt_txt2img",
variant='fp16',
torch_dtype=torch.float16,
scheduler=scheduler,)
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
custom_pipeline="stable_diffusion_tensorrt_txt2img",
variant='fp16',
torch_dtype=torch.float16,
scheduler=scheduler,)
# re-use cached folder to save ONNX models and TensorRT Engines
pipe.set_cached_folder("stabilityai/stable-diffusion-2-1", variant='fp16',)
@@ -1641,18 +1640,18 @@ from io import BytesIO
from PIL import Image
import torch
from diffusers import DDIMScheduler
from diffusers.pipelines.stable_diffusion import StableDiffusionImg2ImgPipeline
from diffusers import DiffusionPipeline
# Use the DDIMScheduler scheduler here instead
scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1",
subfolder="scheduler")
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
custom_pipeline="stable_diffusion_tensorrt_img2img",
variant='fp16',
torch_dtype=torch.float16,
scheduler=scheduler,)
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
custom_pipeline="stable_diffusion_tensorrt_img2img",
variant='fp16',
torch_dtype=torch.float16,
scheduler=scheduler,)
# re-use cached folder to save ONNX models and TensorRT Engines
pipe.set_cached_folder("stabilityai/stable-diffusion-2-1", variant='fp16',)
@@ -2231,12 +2230,12 @@ from io import BytesIO
from PIL import Image
import torch
from diffusers import PNDMScheduler
from diffusers.pipelines.stable_diffusion import StableDiffusionInpaintPipeline
from diffusers.pipelines import DiffusionPipeline
# Use the PNDMScheduler scheduler here instead
scheduler = PNDMScheduler.from_pretrained("stabilityai/stable-diffusion-2-inpainting", subfolder="scheduler")
pipe = StableDiffusionInpaintPipeline.from_pretrained("stabilityai/stable-diffusion-2-inpainting",
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-inpainting",
custom_pipeline="stable_diffusion_tensorrt_inpaint",
variant='fp16',
torch_dtype=torch.float16,
+6 -6
View File
@@ -26,7 +26,7 @@ from gmflow.gmflow import GMFlow
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, ControlNetModel, ImageProjection, UNet2DConditionModel
from diffusers.models.attention_processor import AttnProcessor2_0
from diffusers.models.lora import adjust_lora_scale_text_encoder
@@ -1252,8 +1252,8 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -1456,7 +1456,7 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -1588,7 +1588,7 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -2436,7 +2436,7 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
)
if guess_mode and self.do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
+4 -4
View File
@@ -7,7 +7,7 @@ from transformers import AutoModel, AutoTokenizer, CLIPImageProcessor
from diffusers import DiffusionPipeline
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.pipeline_utils import StableDiffusionMixin
@@ -194,7 +194,7 @@ def retrieve_timesteps(
return timesteps, num_inference_steps
class GlueGenStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin, LoraLoaderMixin):
class GlueGenStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin, StableDiffusionLoraLoaderMixin):
def __init__(
self,
vae: AutoencoderKL,
@@ -290,7 +290,7 @@ class GlueGenStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin, Lo
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -424,7 +424,7 @@ class GlueGenStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin, Lo
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
+9 -5
View File
@@ -21,7 +21,7 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import FromSingleFileMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.pipeline_utils import DiffusionPipeline, StableDiffusionMixin
@@ -53,7 +53,11 @@ def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0):
class InstaFlowPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
StableDiffusionLoraLoaderMixin,
FromSingleFileMixin,
):
r"""
Pipeline for text-to-image generation using Rectified Flow and Euler discretization.
@@ -64,8 +68,8 @@ class InstaFlowPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
@@ -251,7 +255,7 @@ class InstaFlowPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
+11 -6
View File
@@ -24,7 +24,12 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPV
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.attention_processor import (
AttnProcessor,
@@ -130,7 +135,7 @@ class IPAdapterFaceIDStableDiffusionPipeline(
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
IPAdapterMixin,
FromSingleFileMixin,
):
@@ -142,8 +147,8 @@ class IPAdapterFaceIDStableDiffusionPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -518,7 +523,7 @@ class IPAdapterFaceIDStableDiffusionPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -650,7 +655,7 @@ class IPAdapterFaceIDStableDiffusionPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
+2 -2
View File
@@ -395,8 +395,8 @@ class StableDiffusionHighResFixPipeline(StableDiffusionPipeline):
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -6,7 +6,7 @@ import torch
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import FromSingleFileMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.pipeline_utils import DiffusionPipeline, StableDiffusionMixin
@@ -190,7 +190,11 @@ def slerp(
class LatentConsistencyModelWalkPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
StableDiffusionLoraLoaderMixin,
FromSingleFileMixin,
):
r"""
Pipeline for text-to-image generation using a latent consistency model.
@@ -200,8 +204,8 @@ class LatentConsistencyModelWalkPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
@@ -317,7 +321,7 @@ class LatentConsistencyModelWalkPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -449,7 +453,7 @@ class LatentConsistencyModelWalkPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
+9 -4
View File
@@ -29,7 +29,12 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPV
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.attention import Attention, GatedSelfAttentionDense
from diffusers.models.attention_processor import AttnProcessor2_0
@@ -271,7 +276,7 @@ class LLMGroundedDiffusionPipeline(
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
IPAdapterMixin,
FromSingleFileMixin,
):
@@ -1263,7 +1268,7 @@ class LLMGroundedDiffusionPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -1397,7 +1402,7 @@ class LLMGroundedDiffusionPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
+64 -8
View File
@@ -11,15 +11,19 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers import DiffusionPipeline
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import FromSingleFileMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.pipeline_utils import StableDiffusionMixin
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput, StableDiffusionSafetyChecker
from diffusers.schedulers import KarrasDiffusionSchedulers
from diffusers.utils import (
PIL_INTERPOLATION,
USE_PEFT_BACKEND,
deprecate,
logging,
scale_lora_layers,
unscale_lora_layers,
)
from diffusers.utils.torch_utils import randn_tensor
@@ -199,6 +203,7 @@ def get_unweighted_text_embeddings(
text_input: torch.Tensor,
chunk_length: int,
no_boseos_middle: Optional[bool] = True,
clip_skip: Optional[int] = None,
):
"""
When the length of tokens is a multiple of the capacity of the text encoder,
@@ -214,7 +219,20 @@ def get_unweighted_text_embeddings(
# cover the head and the tail by the starting and the ending tokens
text_input_chunk[:, 0] = text_input[0, 0]
text_input_chunk[:, -1] = text_input[0, -1]
text_embedding = pipe.text_encoder(text_input_chunk)[0]
if clip_skip is None:
prompt_embeds = pipe.text_encoder(text_input_chunk.to(pipe.device))
text_embedding = prompt_embeds[0]
else:
prompt_embeds = pipe.text_encoder(text_input_chunk.to(pipe.device), output_hidden_states=True)
# Access the `hidden_states` first, that contains a tuple of
# all the hidden states from the encoder layers. Then index into
# the tuple to access the hidden states from the desired layer.
prompt_embeds = prompt_embeds[-1][-(clip_skip + 1)]
# We also need to apply the final LayerNorm here to not mess with the
# representations. The `last_hidden_states` that we typically use for
# obtaining the final prompt representations passes through the LayerNorm
# layer.
text_embedding = pipe.text_encoder.text_model.final_layer_norm(prompt_embeds)
if no_boseos_middle:
if i == 0:
@@ -230,7 +248,10 @@ def get_unweighted_text_embeddings(
text_embeddings.append(text_embedding)
text_embeddings = torch.concat(text_embeddings, axis=1)
else:
text_embeddings = pipe.text_encoder(text_input)[0]
if clip_skip is None:
clip_skip = 0
prompt_embeds = pipe.text_encoder(text_input, output_hidden_states=True)[-1][-(clip_skip + 1)]
text_embeddings = pipe.text_encoder.text_model.final_layer_norm(prompt_embeds)
return text_embeddings
@@ -242,6 +263,8 @@ def get_weighted_text_embeddings(
no_boseos_middle: Optional[bool] = False,
skip_parsing: Optional[bool] = False,
skip_weighting: Optional[bool] = False,
clip_skip=None,
lora_scale=None,
):
r"""
Prompts can be assigned with local weights using brackets. For example,
@@ -268,6 +291,16 @@ def get_weighted_text_embeddings(
skip_weighting (`bool`, *optional*, defaults to `False`):
Skip the weighting. When the parsing is skipped, it is forced True.
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(pipe, StableDiffusionLoraLoaderMixin):
pipe._lora_scale = lora_scale
# dynamically adjust the LoRA scale
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(pipe.text_encoder, lora_scale)
else:
scale_lora_layers(pipe.text_encoder, lora_scale)
max_length = (pipe.tokenizer.model_max_length - 2) * max_embeddings_multiples + 2
if isinstance(prompt, str):
prompt = [prompt]
@@ -334,10 +367,7 @@ def get_weighted_text_embeddings(
# get the embeddings
text_embeddings = get_unweighted_text_embeddings(
pipe,
prompt_tokens,
pipe.tokenizer.model_max_length,
no_boseos_middle=no_boseos_middle,
pipe, prompt_tokens, pipe.tokenizer.model_max_length, no_boseos_middle=no_boseos_middle, clip_skip=clip_skip
)
prompt_weights = torch.tensor(prompt_weights, dtype=text_embeddings.dtype, device=text_embeddings.device)
if uncond_prompt is not None:
@@ -346,6 +376,7 @@ def get_weighted_text_embeddings(
uncond_tokens,
pipe.tokenizer.model_max_length,
no_boseos_middle=no_boseos_middle,
clip_skip=clip_skip,
)
uncond_weights = torch.tensor(uncond_weights, dtype=uncond_embeddings.dtype, device=uncond_embeddings.device)
@@ -362,6 +393,11 @@ def get_weighted_text_embeddings(
current_mean = uncond_embeddings.float().mean(axis=[-2, -1]).to(uncond_embeddings.dtype)
uncond_embeddings *= (previous_mean / current_mean).unsqueeze(-1).unsqueeze(-1)
if pipe.text_encoder is not None:
if isinstance(pipe, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(pipe.text_encoder, lora_scale)
if uncond_prompt is not None:
return text_embeddings, uncond_embeddings
return text_embeddings, None
@@ -409,7 +445,11 @@ def preprocess_mask(mask, batch_size, scale_factor=8):
class StableDiffusionLongPromptWeightingPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
StableDiffusionLoraLoaderMixin,
FromSingleFileMixin,
):
r"""
Pipeline for text-to-image generation using Stable Diffusion without tokens length limit, and support parsing
@@ -545,6 +585,8 @@ class StableDiffusionLongPromptWeightingPipeline(
max_embeddings_multiples=3,
prompt_embeds: Optional[torch.Tensor] = None,
negative_prompt_embeds: Optional[torch.Tensor] = None,
clip_skip: Optional[int] = None,
lora_scale: Optional[float] = None,
):
r"""
Encodes the prompt into text encoder hidden states.
@@ -593,6 +635,8 @@ class StableDiffusionLongPromptWeightingPipeline(
prompt=prompt,
uncond_prompt=negative_prompt if do_classifier_free_guidance else None,
max_embeddings_multiples=max_embeddings_multiples,
clip_skip=clip_skip,
lora_scale=lora_scale,
)
if prompt_embeds is None:
prompt_embeds = prompt_embeds1
@@ -786,6 +830,7 @@ class StableDiffusionLongPromptWeightingPipeline(
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.Tensor], None]] = None,
is_cancelled_callback: Optional[Callable[[], bool]] = None,
clip_skip: Optional[int] = None,
callback_steps: int = 1,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
):
@@ -861,6 +906,9 @@ class StableDiffusionLongPromptWeightingPipeline(
is_cancelled_callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. If the function returns
`True`, the inference will be cancelled.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function will be called. If not specified, the callback will be
called at every step.
@@ -899,6 +947,7 @@ class StableDiffusionLongPromptWeightingPipeline(
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
lora_scale = cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None
# 3. Encode input prompt
prompt_embeds = self._encode_prompt(
@@ -910,6 +959,8 @@ class StableDiffusionLongPromptWeightingPipeline(
max_embeddings_multiples,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
clip_skip=clip_skip,
lora_scale=lora_scale,
)
dtype = prompt_embeds.dtype
@@ -1040,6 +1091,7 @@ class StableDiffusionLongPromptWeightingPipeline(
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.Tensor], None]] = None,
is_cancelled_callback: Optional[Callable[[], bool]] = None,
clip_skip=None,
callback_steps: int = 1,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
):
@@ -1097,6 +1149,9 @@ class StableDiffusionLongPromptWeightingPipeline(
is_cancelled_callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. If the function returns
`True`, the inference will be cancelled.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function will be called. If not specified, the callback will be
called at every step.
@@ -1131,6 +1186,7 @@ class StableDiffusionLongPromptWeightingPipeline(
return_dict=return_dict,
callback=callback,
is_cancelled_callback=is_cancelled_callback,
clip_skip=clip_skip,
callback_steps=callback_steps,
cross_attention_kwargs=cross_attention_kwargs,
)
+47 -6
View File
@@ -22,19 +22,28 @@ from transformers import (
from diffusers import DiffusionPipeline, StableDiffusionXLPipeline
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionXLLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, ImageProjection, UNet2DConditionModel
from diffusers.models.attention_processor import AttnProcessor2_0, XFormersAttnProcessor
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.pipeline_utils import StableDiffusionMixin
from diffusers.pipelines.stable_diffusion_xl.pipeline_output import StableDiffusionXLPipelineOutput
from diffusers.schedulers import KarrasDiffusionSchedulers
from diffusers.utils import (
USE_PEFT_BACKEND,
deprecate,
is_accelerate_available,
is_accelerate_version,
is_invisible_watermark_available,
logging,
replace_example_docstring,
scale_lora_layers,
unscale_lora_layers,
)
from diffusers.utils.torch_utils import randn_tensor
@@ -256,6 +265,7 @@ def get_weighted_text_embeddings_sdxl(
num_images_per_prompt: int = 1,
device: Optional[torch.device] = None,
clip_skip: Optional[int] = None,
lora_scale: Optional[int] = None,
):
"""
This function can process long prompt with weights, no length limitation
@@ -276,6 +286,24 @@ def get_weighted_text_embeddings_sdxl(
"""
device = device or pipe._execution_device
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(pipe, StableDiffusionXLLoraLoaderMixin):
pipe._lora_scale = lora_scale
# dynamically adjust the LoRA scale
if pipe.text_encoder is not None:
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(pipe.text_encoder, lora_scale)
else:
scale_lora_layers(pipe.text_encoder, lora_scale)
if pipe.text_encoder_2 is not None:
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(pipe.text_encoder_2, lora_scale)
else:
scale_lora_layers(pipe.text_encoder_2, lora_scale)
if prompt_2:
prompt = f"{prompt} {prompt_2}"
@@ -424,6 +452,16 @@ def get_weighted_text_embeddings_sdxl(
bs_embed * num_images_per_prompt, -1
)
if pipe.text_encoder is not None:
if isinstance(pipe, StableDiffusionXLLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(pipe.text_encoder, lora_scale)
if pipe.text_encoder_2 is not None:
if isinstance(pipe, StableDiffusionXLLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(pipe.text_encoder_2, lora_scale)
return prompt_embeds, negative_prompt_embeds, pooled_prompt_embeds, negative_pooled_prompt_embeds
@@ -544,7 +582,7 @@ class SDXLLongPromptWeightingPipeline(
StableDiffusionMixin,
FromSingleFileMixin,
IPAdapterMixin,
LoraLoaderMixin,
StableDiffusionXLLoraLoaderMixin,
TextualInversionLoaderMixin,
):
r"""
@@ -556,8 +594,8 @@ class SDXLLongPromptWeightingPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
Args:
@@ -738,7 +776,7 @@ class SDXLLongPromptWeightingPipeline(
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionXLLoraLoaderMixin):
self._lora_scale = lora_scale
if prompt is not None and isinstance(prompt, str):
@@ -1607,7 +1645,9 @@ class SDXLLongPromptWeightingPipeline(
image_embeds = torch.cat([negative_image_embeds, image_embeds])
# 3. Encode input prompt
(self.cross_attention_kwargs.get("scale", None) if self.cross_attention_kwargs is not None else None)
lora_scale = (
self._cross_attention_kwargs.get("scale", None) if self._cross_attention_kwargs is not None else None
)
negative_prompt = negative_prompt if negative_prompt is not None else ""
@@ -1622,6 +1662,7 @@ class SDXLLongPromptWeightingPipeline(
neg_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
clip_skip=clip_skip,
lora_scale=lora_scale,
)
dtype = prompt_embeds.dtype
@@ -22,7 +22,7 @@ from PIL import Image
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import IPAdapterMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, ControlNetModel, ImageProjection, UNet2DConditionModel, UNetMotionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.models.unets.unet_motion_model import MotionAdapter
@@ -114,7 +114,11 @@ def tensor2vid(video: torch.Tensor, processor, output_type="np"):
class AnimateDiffControlNetPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, IPAdapterMixin, LoraLoaderMixin
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
):
r"""
Pipeline for text-to-video generation.
@@ -124,8 +128,8 @@ class AnimateDiffControlNetPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
Args:
@@ -234,7 +238,7 @@ class AnimateDiffControlNetPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -366,7 +370,7 @@ class AnimateDiffControlNetPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -27,7 +27,7 @@ import torch
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import IPAdapterMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, ImageProjection, UNet2DConditionModel, UNetMotionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.models.unet_motion_model import MotionAdapter
@@ -240,7 +240,11 @@ def retrieve_timesteps(
class AnimateDiffImgToVideoPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, IPAdapterMixin, LoraLoaderMixin
DiffusionPipeline,
StableDiffusionMixin,
TextualInversionLoaderMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
):
r"""
Pipeline for image-to-video generation.
@@ -250,8 +254,8 @@ class AnimateDiffImgToVideoPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
Args:
@@ -351,7 +355,7 @@ class AnimateDiffImgToVideoPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -483,7 +487,7 @@ class AnimateDiffImgToVideoPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -12,7 +12,7 @@ from transformers import CLIPTextModel, CLIPTextModelWithProjection, CLIPTokeniz
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import (
FromSingleFileMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, UNet2DConditionModel
@@ -89,7 +89,11 @@ def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0):
class DemoFusionSDXLPipeline(
DiffusionPipeline, StableDiffusionMixin, FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
DiffusionPipeline,
StableDiffusionMixin,
FromSingleFileMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
):
r"""
Pipeline for text-to-image generation using Stable Diffusion XL.
@@ -231,7 +235,7 @@ class DemoFusionSDXLPipeline(
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
+2 -2
View File
@@ -21,7 +21,7 @@ from transformers import CLIPTextModel, CLIPTokenizer
from diffusers import AutoencoderKL, UNet2DConditionModel
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models.attention import BasicTransformerBlock
from diffusers.models.attention_processor import LoRAAttnProcessor
from diffusers.pipelines.pipeline_utils import DiffusionPipeline
@@ -222,7 +222,7 @@ class FabricPipeline(DiffusionPipeline):
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
if prompt is not None and isinstance(prompt, str):
+6 -6
View File
@@ -35,7 +35,7 @@ from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models.attention import Attention
@@ -75,7 +75,7 @@ def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0):
class Prompt2PromptPipeline(
DiffusionPipeline,
TextualInversionLoaderMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
IPAdapterMixin,
FromSingleFileMixin,
):
@@ -87,8 +87,8 @@ class Prompt2PromptPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -286,7 +286,7 @@ class Prompt2PromptPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -420,7 +420,7 @@ class Prompt2PromptPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -27,7 +27,12 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPV
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, ImageProjection, UNet2DConditionModel
from diffusers.models.attention_processor import Attention, FusedAttnProcessor2_0
from diffusers.models.lora import adjust_lora_scale_text_encoder
@@ -358,7 +363,7 @@ def retrieve_timesteps(
class StableDiffusionBoxDiffPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
DiffusionPipeline, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion with BoxDiff.
@@ -368,8 +373,8 @@ class StableDiffusionBoxDiffPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -594,7 +599,7 @@ class StableDiffusionBoxDiffPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -726,7 +731,7 @@ class StableDiffusionBoxDiffPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -11,7 +11,12 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPV
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models import AutoencoderKL, ImageProjection, UNet2DConditionModel
from diffusers.models.attention_processor import Attention, AttnProcessor2_0, FusedAttnProcessor2_0
from diffusers.models.lora import adjust_lora_scale_text_encoder
@@ -328,7 +333,7 @@ def retrieve_timesteps(
class StableDiffusionPAGPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
DiffusionPipeline, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion.
@@ -336,8 +341,8 @@ class StableDiffusionPAGPipeline(
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
Args:
@@ -560,7 +565,7 @@ class StableDiffusionPAGPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -692,7 +697,7 @@ class StableDiffusionPAGPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -22,7 +22,7 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers import DiffusionPipeline
from diffusers.image_processor import PipelineDepthInput, PipelineImageInput, VaeImageProcessorLDM3D
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import FromSingleFileMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
@@ -69,7 +69,7 @@ EXAMPLE_DOC_STRING = """
class StableDiffusionUpscaleLDM3DPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
DiffusionPipeline, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image and 3D generation using LDM3D.
@@ -79,8 +79,8 @@ class StableDiffusionUpscaleLDM3DPipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
@@ -233,7 +233,7 @@ class StableDiffusionUpscaleLDM3DPipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -365,7 +365,7 @@ class StableDiffusionUpscaleLDM3DPipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -33,7 +33,7 @@ from diffusers import DiffusionPipeline
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import (
FromSingleFileMixin,
LoraLoaderMixin,
StableDiffusionLoraLoaderMixin,
StableDiffusionXLLoraLoaderMixin,
TextualInversionLoaderMixin,
)
@@ -300,7 +300,7 @@ def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0):
class StableDiffusionXLControlNetAdapterInpaintPipeline(
DiffusionPipeline, StableDiffusionMixin, FromSingleFileMixin, LoraLoaderMixin
DiffusionPipeline, StableDiffusionMixin, FromSingleFileMixin, StableDiffusionLoraLoaderMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion augmented with T2I-Adapter
@@ -178,11 +178,11 @@ class StableDiffusionXLDifferentialImg2ImgPipeline(
In addition the pipeline inherits the following loading methods:
- *Textual-Inversion*: [`loaders.TextualInversionLoaderMixin.load_textual_inversion`]
- *LoRA*: [`loaders.LoraLoaderMixin.load_lora_weights`]
- *LoRA*: [`loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`]
- *Ckpt*: [`loaders.FromSingleFileMixin.from_single_file`]
as well as the following saving methods:
- *LoRA*: [`loaders.LoraLoaderMixin.save_lora_weights`]
- *LoRA*: [`loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`]
Args:
vae ([`AutoencoderKL`]):
@@ -1002,7 +1002,7 @@ class StableDiffusionXLInstantIDImg2ImgPipeline(StableDiffusionXLControlNetImg2I
)
if guess_mode and self.do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
@@ -991,7 +991,7 @@ class StableDiffusionXLInstantIDPipeline(StableDiffusionXLControlNetPipeline):
)
if guess_mode and self.do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
+2 -2
View File
@@ -864,7 +864,7 @@ class RerenderAVideoPipeline(StableDiffusionControlNetImg2ImgPipeline):
)
if guess_mode and do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
@@ -1038,7 +1038,7 @@ class RerenderAVideoPipeline(StableDiffusionControlNetImg2ImgPipeline):
)
if guess_mode and do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [
+2 -2
View File
@@ -11,7 +11,7 @@ from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
from diffusers import AutoencoderKL, DiffusionPipeline, DPMSolverMultistepScheduler, UNet2DConditionModel
from diffusers.loaders import AttnProcsLayers, LoraLoaderMixin
from diffusers.loaders import AttnProcsLayers, StableDiffusionLoraLoaderMixin
from diffusers.models.attention_processor import (
AttnAddedKVProcessor,
AttnAddedKVProcessor2_0,
@@ -321,7 +321,7 @@ class SdeDragPipeline(DiffusionPipeline):
optimizer.zero_grad()
with tempfile.TemporaryDirectory() as save_lora_dir:
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
save_directory=save_lora_dir,
unet_lora_layers=unet_lora_layers,
text_encoder_lora_layers=None,
@@ -752,7 +752,7 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
)
if guess_mode and do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
+2 -2
View File
@@ -21,7 +21,7 @@ from packaging import version
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from diffusers.configuration_utils import FrozenDict
from diffusers.loaders import LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.pipeline_utils import DiffusionPipeline, StableDiffusionMixin
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
@@ -61,7 +61,7 @@ EXAMPLE_DOC_STRING = """
class StableDiffusionIPEXPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, LoraLoaderMixin
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion on IPEX.
@@ -11,7 +11,12 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers import AutoencoderKL, DiffusionPipeline, UNet2DConditionModel
from diffusers.configuration_utils import FrozenDict, deprecate
from diffusers.image_processor import VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import (
FromSingleFileMixin,
IPAdapterMixin,
StableDiffusionLoraLoaderMixin,
TextualInversionLoaderMixin,
)
from diffusers.models.attention import BasicTransformerBlock
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.models.unets.unet_2d_blocks import CrossAttnDownBlock2D, CrossAttnUpBlock2D, DownBlock2D, UpBlock2D
@@ -76,7 +81,7 @@ def torch_dfs(model: torch.nn.Module):
class StableDiffusionReferencePipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
DiffusionPipeline, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
):
r"""
Pipeline for Stable Diffusion Reference.
@@ -86,8 +91,8 @@ class StableDiffusionReferencePipeline(
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters
@@ -443,7 +448,7 @@ class StableDiffusionReferencePipeline(
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -575,7 +580,7 @@ class StableDiffusionReferencePipeline(
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -23,7 +23,7 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers import AutoencoderKL, DiffusionPipeline, UNet2DConditionModel
from diffusers.configuration_utils import FrozenDict, deprecate
from diffusers.loaders import LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.pipelines.pipeline_utils import StableDiffusionMixin
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.pipelines.stable_diffusion.safety_checker import (
@@ -140,7 +140,7 @@ def prepare_mask_and_masked_image(image, mask):
class StableDiffusionRepaintPipeline(
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, LoraLoaderMixin
DiffusionPipeline, StableDiffusionMixin, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin
):
r"""
Pipeline for text-guided image inpainting using Stable Diffusion. *This is an experimental feature*.
@@ -148,9 +148,9 @@ class StableDiffusionRepaintPipeline(
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
In addition the pipeline inherits the following loading methods:
- *Textual-Inversion*: [`loaders.TextualInversionLoaderMixin.load_textual_inversion`]
- *LoRA*: [`loaders.LoraLoaderMixin.load_lora_weights`]
- *LoRA*: [`loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`]
as well as the following saving methods:
- *LoRA*: [`loaders.LoraLoaderMixin.save_lora_weights`]
- *LoRA*: [`loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`]
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
@@ -18,8 +18,7 @@
import gc
import os
from collections import OrderedDict
from copy import copy
from typing import List, Optional, Union
from typing import List, Optional, Tuple, Union
import numpy as np
import onnx
@@ -27,9 +26,11 @@ import onnx_graphsurgeon as gs
import PIL.Image
import tensorrt as trt
import torch
from cuda import cudart
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from packaging import version
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
from polygraphy.backend.onnx.loader import fold_constants
@@ -41,12 +42,13 @@ from polygraphy.backend.trt import (
network_from_onnx_path,
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers import DiffusionPipeline
from diffusers.configuration_utils import FrozenDict, deprecate
from diffusers.image_processor import VaeImageProcessor
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
StableDiffusionImg2ImgPipeline,
StableDiffusionPipelineOutput,
StableDiffusionSafetyChecker,
)
@@ -58,7 +60,7 @@ from diffusers.utils import logging
"""
Installation instructions
python3 -m pip install --upgrade transformers diffusers>=0.16.0
python3 -m pip install --upgrade tensorrt>=8.6.1
python3 -m pip install --upgrade tensorrt~=10.2.0
python3 -m pip install --upgrade polygraphy>=0.47.0 onnx-graphsurgeon --extra-index-url https://pypi.ngc.nvidia.com
python3 -m pip install onnxruntime
"""
@@ -88,10 +90,6 @@ else:
torch_to_numpy_dtype_dict = {value: key for (key, value) in numpy_to_torch_dtype_dict.items()}
def device_view(t):
return cuda.DeviceView(ptr=t.data_ptr(), shape=t.shape, dtype=torch_to_numpy_dtype_dict[t.dtype])
def preprocess_image(image):
"""
image: torch.Tensor
@@ -125,10 +123,8 @@ class Engine:
onnx_path,
fp16,
input_profile=None,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
workspace_size=0,
):
logger.warning(f"Building TensorRT engine for {onnx_path}: {self.engine_path}")
p = Profile()
@@ -137,20 +133,13 @@ class Engine:
assert len(dims) == 3
p.add(name, min=dims[0], opt=dims[1], max=dims[2])
config_kwargs = {}
config_kwargs["preview_features"] = [trt.PreviewFeature.DISABLE_EXTERNAL_TACTIC_SOURCES_FOR_CORE_0805]
if enable_preview:
# Faster dynamic shapes made optional since it increases engine build time.
config_kwargs["preview_features"].append(trt.PreviewFeature.FASTER_DYNAMIC_SHAPES_0805)
if workspace_size > 0:
config_kwargs["memory_pool_limits"] = {trt.MemoryPoolType.WORKSPACE: workspace_size}
extra_build_args = {}
if not enable_all_tactics:
config_kwargs["tactic_sources"] = []
extra_build_args["tactic_sources"] = []
engine = engine_from_network(
network_from_onnx_path(onnx_path, flags=[trt.OnnxParserFlag.NATIVE_INSTANCENORM]),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **config_kwargs),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **extra_build_args),
save_timing_cache=timing_cache,
)
save_engine(engine, path=self.engine_path)
@@ -163,28 +152,24 @@ class Engine:
self.context = self.engine.create_execution_context()
def allocate_buffers(self, shape_dict=None, device="cuda"):
for idx in range(trt_util.get_bindings_per_profile(self.engine)):
binding = self.engine[idx]
if shape_dict and binding in shape_dict:
shape = shape_dict[binding]
for binding in range(self.engine.num_io_tensors):
name = self.engine.get_tensor_name(binding)
if shape_dict and name in shape_dict:
shape = shape_dict[name]
else:
shape = self.engine.get_binding_shape(binding)
dtype = trt.nptype(self.engine.get_binding_dtype(binding))
if self.engine.binding_is_input(binding):
self.context.set_binding_shape(idx, shape)
shape = self.engine.get_tensor_shape(name)
dtype = trt.nptype(self.engine.get_tensor_dtype(name))
if self.engine.get_tensor_mode(name) == trt.TensorIOMode.INPUT:
self.context.set_input_shape(name, shape)
tensor = torch.empty(tuple(shape), dtype=numpy_to_torch_dtype_dict[dtype]).to(device=device)
self.tensors[binding] = tensor
self.buffers[binding] = cuda.DeviceView(ptr=tensor.data_ptr(), shape=shape, dtype=dtype)
self.tensors[name] = tensor
def infer(self, feed_dict, stream):
start_binding, end_binding = trt_util.get_active_profile_bindings(self.context)
# shallow copy of ordered dict
device_buffers = copy(self.buffers)
for name, buf in feed_dict.items():
assert isinstance(buf, cuda.DeviceView)
device_buffers[name] = buf
bindings = [0] * start_binding + [buf.ptr for buf in device_buffers.values()]
noerror = self.context.execute_async_v2(bindings=bindings, stream_handle=stream.ptr)
self.tensors[name].copy_(buf)
for name, tensor in self.tensors.items():
self.context.set_tensor_address(name, tensor.data_ptr())
noerror = self.context.execute_async_v3(stream)
if not noerror:
raise ValueError("ERROR: inference failed.")
@@ -325,10 +310,8 @@ def build_engines(
force_engine_rebuild=False,
static_batch=False,
static_shape=True,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
max_workspace_size=0,
):
built_engines = {}
if not os.path.isdir(onnx_dir):
@@ -393,9 +376,7 @@ def build_engines(
static_batch=static_batch,
static_shape=static_shape,
),
enable_preview=enable_preview,
timing_cache=timing_cache,
workspace_size=max_workspace_size,
)
built_engines[model_name] = engine
@@ -674,11 +655,11 @@ def make_VAEEncoder(model, device, max_batch_size, embedding_dim, inpaint=False)
return VAEEncoder(model, device=device, max_batch_size=max_batch_size, embedding_dim=embedding_dim)
class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
class TensorRTStableDiffusionImg2ImgPipeline(DiffusionPipeline):
r"""
Pipeline for image-to-image generation using TensorRT accelerated Stable Diffusion.
This model inherits from [`StableDiffusionImg2ImgPipeline`]. Check the superclass documentation for the generic methods the
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
Args:
@@ -702,6 +683,8 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor", "image_encoder"]
def __init__(
self,
vae: AutoencoderKL,
@@ -722,24 +705,86 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
onnx_dir: str = "onnx",
# TensorRT engine build parameters
engine_dir: str = "engine",
build_preview_features: bool = True,
force_engine_rebuild: bool = False,
timing_cache: str = "timing_cache",
):
super().__init__(
vae,
text_encoder,
tokenizer,
unet,
scheduler,
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} is outdated. `steps_offset`"
f" should be set to 1 instead of {scheduler.config.steps_offset}. Please make sure "
"to update the config accordingly as leaving `steps_offset` might led to incorrect results"
" in future versions. If you have downloaded this checkpoint from the Hugging Face Hub,"
" it would be very nice if you could open a Pull request for the `scheduler/scheduler_config.json`"
" file"
)
deprecate("steps_offset!=1", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["steps_offset"] = 1
scheduler._internal_dict = FrozenDict(new_config)
if hasattr(scheduler.config, "clip_sample") and scheduler.config.clip_sample is True:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} has not set the configuration `clip_sample`."
" `clip_sample` should be set to False in the configuration file. Please make sure to update the"
" config accordingly as not setting `clip_sample` in the config might lead to incorrect results in"
" future versions. If you have downloaded this checkpoint from the Hugging Face Hub, it would be very"
" nice if you could open a Pull request for the `scheduler/scheduler_config.json` file"
)
deprecate("clip_sample not set", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["clip_sample"] = False
scheduler._internal_dict = FrozenDict(new_config)
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
version.parse(unet.config._diffusers_version).base_version
) < version.parse("0.9.0.dev0")
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
deprecation_message = (
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
" the `unet/config.json` file"
)
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(unet.config)
new_config["sample_size"] = 64
unet._internal_dict = FrozenDict(new_config)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
self.stages = stages
self.image_height, self.image_width = image_height, image_width
self.inpaint = False
@@ -750,7 +795,6 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
self.timing_cache = timing_cache
self.build_static_batch = False
self.build_dynamic_shape = False
self.build_preview_features = build_preview_features
self.max_batch_size = max_batch_size
# TODO: Restrict batch size to 4 for larger image dimensions as a WAR for TensorRT limitation.
@@ -761,6 +805,11 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
self.models = {} # loaded in __loadModels()
self.engine = {} # loaded in build_engines()
self.vae.forward = self.vae.decode
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor)
self.register_to_config(requires_safety_checker=requires_safety_checker)
def __loadModels(self):
# Load pipeline models
self.embedding_dim = self.text_encoder.config.hidden_size
@@ -779,6 +828,33 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
if "vae_encoder" in self.stages:
self.models["vae_encoder"] = make_VAEEncoder(self.vae, **models_args)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker
def run_safety_checker(
self, image: Union[torch.Tensor, PIL.Image.Image], device: torch.device, dtype: torch.dtype
) -> Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]:
r"""
Runs the safety checker on the given image.
Args:
image (Union[torch.Tensor, PIL.Image.Image]): The input image to be checked.
device (torch.device): The device to run the safety checker on.
dtype (torch.dtype): The data type of the input image.
Returns:
(image, has_nsfw_concept) Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]: A tuple containing the processed image and
a boolean indicating whether the image has a NSFW (Not Safe for Work) concept.
"""
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
@@ -826,7 +902,6 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
force_engine_rebuild=self.force_engine_rebuild,
static_batch=self.build_static_batch,
static_shape=not self.build_dynamic_shape,
enable_preview=self.build_preview_features,
timing_cache=self.timing_cache,
)
@@ -850,9 +925,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
return tuple(init_images)
def __encode_image(self, init_image):
init_latents = runEngine(self.engine["vae_encoder"], {"images": device_view(init_image)}, self.stream)[
"latent"
]
init_latents = runEngine(self.engine["vae_encoder"], {"images": init_image}, self.stream)["latent"]
init_latents = 0.18215 * init_latents
return init_latents
@@ -881,9 +954,8 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
.to(self.torch_device)
)
text_input_ids_inp = device_view(text_input_ids)
# NOTE: output tensor for CLIP must be cloned because it will be overwritten when called again for negative prompt
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids_inp}, self.stream)[
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids}, self.stream)[
"text_embeddings"
].clone()
@@ -899,8 +971,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
.input_ids.type(torch.int32)
.to(self.torch_device)
)
uncond_input_ids_inp = device_view(uncond_input_ids)
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids_inp}, self.stream)[
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids}, self.stream)[
"text_embeddings"
]
@@ -924,18 +995,15 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
# Predict the noise residual
timestep_float = timestep.float() if timestep.dtype != torch.float32 else timestep
sample_inp = device_view(latent_model_input)
timestep_inp = device_view(timestep_float)
embeddings_inp = device_view(text_embeddings)
noise_pred = runEngine(
self.engine["unet"],
{"sample": sample_inp, "timestep": timestep_inp, "encoder_hidden_states": embeddings_inp},
{"sample": latent_model_input, "timestep": timestep_float, "encoder_hidden_states": text_embeddings},
self.stream,
)["latent"]
# Perform guidance
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + self.guidance_scale * (noise_pred_text - noise_pred_uncond)
noise_pred = noise_pred_uncond + self._guidance_scale * (noise_pred_text - noise_pred_uncond)
latents = self.scheduler.step(noise_pred, timestep, latents).prev_sample
@@ -943,12 +1011,12 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
return latents
def __decode_latent(self, latents):
images = runEngine(self.engine["vae"], {"latent": device_view(latents)}, self.stream)["images"]
images = runEngine(self.engine["vae"], {"latent": latents}, self.stream)["images"]
images = (images / 2 + 0.5).clamp(0, 1)
return images.cpu().permute(0, 2, 3, 1).float().numpy()
def __loadResources(self, image_height, image_width, batch_size):
self.stream = cuda.Stream()
self.stream = cudart.cudaStreamCreate()[1]
# Allocate buffers for TensorRT engine bindings
for model_name, obj in self.models.items():
@@ -1061,5 +1129,6 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
# VAE decode latent
images = self.__decode_latent(latents)
images, has_nsfw_concept = self.run_safety_checker(images, self.torch_device, text_embeddings.dtype)
images = self.numpy_to_pil(images)
return StableDiffusionPipelineOutput(images=images, nsfw_content_detected=None)
return StableDiffusionPipelineOutput(images=images, nsfw_content_detected=has_nsfw_concept)
@@ -18,8 +18,7 @@
import gc
import os
from collections import OrderedDict
from copy import copy
from typing import List, Optional, Union
from typing import List, Optional, Tuple, Union
import numpy as np
import onnx
@@ -27,9 +26,11 @@ import onnx_graphsurgeon as gs
import PIL.Image
import tensorrt as trt
import torch
from cuda import cudart
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from packaging import version
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
from polygraphy.backend.onnx.loader import fold_constants
@@ -41,24 +42,29 @@ from polygraphy.backend.trt import (
network_from_onnx_path,
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers import DiffusionPipeline
from diffusers.configuration_utils import FrozenDict, deprecate
from diffusers.image_processor import VaeImageProcessor
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
StableDiffusionInpaintPipeline,
StableDiffusionPipelineOutput,
StableDiffusionSafetyChecker,
)
from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint import prepare_mask_and_masked_image
from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint import (
prepare_mask_and_masked_image,
retrieve_latents,
)
from diffusers.schedulers import DDIMScheduler
from diffusers.utils import logging
from diffusers.utils.torch_utils import randn_tensor
"""
Installation instructions
python3 -m pip install --upgrade transformers diffusers>=0.16.0
python3 -m pip install --upgrade tensorrt>=8.6.1
python3 -m pip install --upgrade tensorrt~=10.2.0
python3 -m pip install --upgrade polygraphy>=0.47.0 onnx-graphsurgeon --extra-index-url https://pypi.ngc.nvidia.com
python3 -m pip install onnxruntime
"""
@@ -88,10 +94,6 @@ else:
torch_to_numpy_dtype_dict = {value: key for (key, value) in numpy_to_torch_dtype_dict.items()}
def device_view(t):
return cuda.DeviceView(ptr=t.data_ptr(), shape=t.shape, dtype=torch_to_numpy_dtype_dict[t.dtype])
def preprocess_image(image):
"""
image: torch.Tensor
@@ -125,10 +127,8 @@ class Engine:
onnx_path,
fp16,
input_profile=None,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
workspace_size=0,
):
logger.warning(f"Building TensorRT engine for {onnx_path}: {self.engine_path}")
p = Profile()
@@ -137,20 +137,13 @@ class Engine:
assert len(dims) == 3
p.add(name, min=dims[0], opt=dims[1], max=dims[2])
config_kwargs = {}
config_kwargs["preview_features"] = [trt.PreviewFeature.DISABLE_EXTERNAL_TACTIC_SOURCES_FOR_CORE_0805]
if enable_preview:
# Faster dynamic shapes made optional since it increases engine build time.
config_kwargs["preview_features"].append(trt.PreviewFeature.FASTER_DYNAMIC_SHAPES_0805)
if workspace_size > 0:
config_kwargs["memory_pool_limits"] = {trt.MemoryPoolType.WORKSPACE: workspace_size}
extra_build_args = {}
if not enable_all_tactics:
config_kwargs["tactic_sources"] = []
extra_build_args["tactic_sources"] = []
engine = engine_from_network(
network_from_onnx_path(onnx_path, flags=[trt.OnnxParserFlag.NATIVE_INSTANCENORM]),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **config_kwargs),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **extra_build_args),
save_timing_cache=timing_cache,
)
save_engine(engine, path=self.engine_path)
@@ -163,28 +156,24 @@ class Engine:
self.context = self.engine.create_execution_context()
def allocate_buffers(self, shape_dict=None, device="cuda"):
for idx in range(trt_util.get_bindings_per_profile(self.engine)):
binding = self.engine[idx]
if shape_dict and binding in shape_dict:
shape = shape_dict[binding]
for binding in range(self.engine.num_io_tensors):
name = self.engine.get_tensor_name(binding)
if shape_dict and name in shape_dict:
shape = shape_dict[name]
else:
shape = self.engine.get_binding_shape(binding)
dtype = trt.nptype(self.engine.get_binding_dtype(binding))
if self.engine.binding_is_input(binding):
self.context.set_binding_shape(idx, shape)
shape = self.engine.get_tensor_shape(name)
dtype = trt.nptype(self.engine.get_tensor_dtype(name))
if self.engine.get_tensor_mode(name) == trt.TensorIOMode.INPUT:
self.context.set_input_shape(name, shape)
tensor = torch.empty(tuple(shape), dtype=numpy_to_torch_dtype_dict[dtype]).to(device=device)
self.tensors[binding] = tensor
self.buffers[binding] = cuda.DeviceView(ptr=tensor.data_ptr(), shape=shape, dtype=dtype)
self.tensors[name] = tensor
def infer(self, feed_dict, stream):
start_binding, end_binding = trt_util.get_active_profile_bindings(self.context)
# shallow copy of ordered dict
device_buffers = copy(self.buffers)
for name, buf in feed_dict.items():
assert isinstance(buf, cuda.DeviceView)
device_buffers[name] = buf
bindings = [0] * start_binding + [buf.ptr for buf in device_buffers.values()]
noerror = self.context.execute_async_v2(bindings=bindings, stream_handle=stream.ptr)
self.tensors[name].copy_(buf)
for name, tensor in self.tensors.items():
self.context.set_tensor_address(name, tensor.data_ptr())
noerror = self.context.execute_async_v3(stream)
if not noerror:
raise ValueError("ERROR: inference failed.")
@@ -325,10 +314,8 @@ def build_engines(
force_engine_rebuild=False,
static_batch=False,
static_shape=True,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
max_workspace_size=0,
):
built_engines = {}
if not os.path.isdir(onnx_dir):
@@ -393,9 +380,7 @@ def build_engines(
static_batch=static_batch,
static_shape=static_shape,
),
enable_preview=enable_preview,
timing_cache=timing_cache,
workspace_size=max_workspace_size,
)
built_engines[model_name] = engine
@@ -674,11 +659,11 @@ def make_VAEEncoder(model, device, max_batch_size, embedding_dim, inpaint=False)
return VAEEncoder(model, device=device, max_batch_size=max_batch_size, embedding_dim=embedding_dim)
class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
class TensorRTStableDiffusionInpaintPipeline(DiffusionPipeline):
r"""
Pipeline for inpainting using TensorRT accelerated Stable Diffusion.
This model inherits from [`StableDiffusionInpaintPipeline`]. Check the superclass documentation for the generic methods the
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
Args:
@@ -702,6 +687,8 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor", "image_encoder"]
def __init__(
self,
vae: AutoencoderKL,
@@ -722,24 +709,86 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
onnx_dir: str = "onnx",
# TensorRT engine build parameters
engine_dir: str = "engine",
build_preview_features: bool = True,
force_engine_rebuild: bool = False,
timing_cache: str = "timing_cache",
):
super().__init__(
vae,
text_encoder,
tokenizer,
unet,
scheduler,
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} is outdated. `steps_offset`"
f" should be set to 1 instead of {scheduler.config.steps_offset}. Please make sure "
"to update the config accordingly as leaving `steps_offset` might led to incorrect results"
" in future versions. If you have downloaded this checkpoint from the Hugging Face Hub,"
" it would be very nice if you could open a Pull request for the `scheduler/scheduler_config.json`"
" file"
)
deprecate("steps_offset!=1", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["steps_offset"] = 1
scheduler._internal_dict = FrozenDict(new_config)
if hasattr(scheduler.config, "clip_sample") and scheduler.config.clip_sample is True:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} has not set the configuration `clip_sample`."
" `clip_sample` should be set to False in the configuration file. Please make sure to update the"
" config accordingly as not setting `clip_sample` in the config might lead to incorrect results in"
" future versions. If you have downloaded this checkpoint from the Hugging Face Hub, it would be very"
" nice if you could open a Pull request for the `scheduler/scheduler_config.json` file"
)
deprecate("clip_sample not set", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["clip_sample"] = False
scheduler._internal_dict = FrozenDict(new_config)
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
version.parse(unet.config._diffusers_version).base_version
) < version.parse("0.9.0.dev0")
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
deprecation_message = (
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
" the `unet/config.json` file"
)
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(unet.config)
new_config["sample_size"] = 64
unet._internal_dict = FrozenDict(new_config)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
self.stages = stages
self.image_height, self.image_width = image_height, image_width
self.inpaint = True
@@ -750,7 +799,6 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
self.timing_cache = timing_cache
self.build_static_batch = False
self.build_dynamic_shape = False
self.build_preview_features = build_preview_features
self.max_batch_size = max_batch_size
# TODO: Restrict batch size to 4 for larger image dimensions as a WAR for TensorRT limitation.
@@ -761,6 +809,11 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
self.models = {} # loaded in __loadModels()
self.engine = {} # loaded in build_engines()
self.vae.forward = self.vae.decode
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor)
self.register_to_config(requires_safety_checker=requires_safety_checker)
def __loadModels(self):
# Load pipeline models
self.embedding_dim = self.text_encoder.config.hidden_size
@@ -779,6 +832,112 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
if "vae_encoder" in self.stages:
self.models["vae_encoder"] = make_VAEEncoder(self.vae, **models_args)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint.StableDiffusionInpaintPipeline
def _encode_vae_image(self, image: torch.Tensor, generator: torch.Generator):
if isinstance(generator, list):
image_latents = [
retrieve_latents(self.vae.encode(image[i : i + 1]), generator=generator[i])
for i in range(image.shape[0])
]
image_latents = torch.cat(image_latents, dim=0)
else:
image_latents = retrieve_latents(self.vae.encode(image), generator=generator)
image_latents = self.vae.config.scaling_factor * image_latents
return image_latents
def prepare_latents(
self,
batch_size,
num_channels_latents,
height,
width,
dtype,
device,
generator,
latents=None,
image=None,
timestep=None,
is_strength_max=True,
return_noise=False,
return_image_latents=False,
):
shape = (
batch_size,
num_channels_latents,
int(height) // self.vae_scale_factor,
int(width) // self.vae_scale_factor,
)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
if (image is None or timestep is None) and not is_strength_max:
raise ValueError(
"Since strength < 1. initial latents are to be initialised as a combination of Image + Noise."
"However, either the image or the noise timestep has not been provided."
)
if return_image_latents or (latents is None and not is_strength_max):
image = image.to(device=device, dtype=dtype)
if image.shape[1] == 4:
image_latents = image
else:
image_latents = self._encode_vae_image(image=image, generator=generator)
image_latents = image_latents.repeat(batch_size // image_latents.shape[0], 1, 1, 1)
if latents is None:
noise = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
# if strength is 1. then initialise the latents to noise, else initial to image + noise
latents = noise if is_strength_max else self.scheduler.add_noise(image_latents, noise, timestep)
# if pure noise then scale the initial latents by the Scheduler's init sigma
latents = latents * self.scheduler.init_noise_sigma if is_strength_max else latents
else:
noise = latents.to(device)
latents = noise * self.scheduler.init_noise_sigma
outputs = (latents,)
if return_noise:
outputs += (noise,)
if return_image_latents:
outputs += (image_latents,)
return outputs
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker
def run_safety_checker(
self, image: Union[torch.Tensor, PIL.Image.Image], device: torch.device, dtype: torch.dtype
) -> Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]:
r"""
Runs the safety checker on the given image.
Args:
image (Union[torch.Tensor, PIL.Image.Image]): The input image to be checked.
device (torch.device): The device to run the safety checker on.
dtype (torch.dtype): The data type of the input image.
Returns:
(image, has_nsfw_concept) Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]: A tuple containing the processed image and
a boolean indicating whether the image has a NSFW (Not Safe for Work) concept.
"""
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
@@ -826,7 +985,6 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
force_engine_rebuild=self.force_engine_rebuild,
static_batch=self.build_static_batch,
static_shape=not self.build_dynamic_shape,
enable_preview=self.build_preview_features,
timing_cache=self.timing_cache,
)
@@ -850,9 +1008,7 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
return tuple(init_images)
def __encode_image(self, init_image):
init_latents = runEngine(self.engine["vae_encoder"], {"images": device_view(init_image)}, self.stream)[
"latent"
]
init_latents = runEngine(self.engine["vae_encoder"], {"images": init_image}, self.stream)["latent"]
init_latents = 0.18215 * init_latents
return init_latents
@@ -881,9 +1037,8 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
.to(self.torch_device)
)
text_input_ids_inp = device_view(text_input_ids)
# NOTE: output tensor for CLIP must be cloned because it will be overwritten when called again for negative prompt
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids_inp}, self.stream)[
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids}, self.stream)[
"text_embeddings"
].clone()
@@ -899,8 +1054,7 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
.input_ids.type(torch.int32)
.to(self.torch_device)
)
uncond_input_ids_inp = device_view(uncond_input_ids)
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids_inp}, self.stream)[
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids}, self.stream)[
"text_embeddings"
]
@@ -924,18 +1078,15 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
# Predict the noise residual
timestep_float = timestep.float() if timestep.dtype != torch.float32 else timestep
sample_inp = device_view(latent_model_input)
timestep_inp = device_view(timestep_float)
embeddings_inp = device_view(text_embeddings)
noise_pred = runEngine(
self.engine["unet"],
{"sample": sample_inp, "timestep": timestep_inp, "encoder_hidden_states": embeddings_inp},
{"sample": latent_model_input, "timestep": timestep_float, "encoder_hidden_states": text_embeddings},
self.stream,
)["latent"]
# Perform guidance
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + self.guidance_scale * (noise_pred_text - noise_pred_uncond)
noise_pred = noise_pred_uncond + self._guidance_scale * (noise_pred_text - noise_pred_uncond)
latents = self.scheduler.step(noise_pred, timestep, latents).prev_sample
@@ -943,12 +1094,12 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
return latents
def __decode_latent(self, latents):
images = runEngine(self.engine["vae"], {"latent": device_view(latents)}, self.stream)["images"]
images = runEngine(self.engine["vae"], {"latent": latents}, self.stream)["images"]
images = (images / 2 + 0.5).clamp(0, 1)
return images.cpu().permute(0, 2, 3, 1).float().numpy()
def __loadResources(self, image_height, image_width, batch_size):
self.stream = cuda.Stream()
self.stream = cudart.cudaStreamCreate()[1]
# Allocate buffers for TensorRT engine bindings
for model_name, obj in self.models.items():
@@ -1112,5 +1263,6 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
# VAE decode latent
images = self.__decode_latent(latents)
images, has_nsfw_concept = self.run_safety_checker(images, self.torch_device, text_embeddings.dtype)
images = self.numpy_to_pil(images)
return StableDiffusionPipelineOutput(images=images, nsfw_content_detected=None)
return StableDiffusionPipelineOutput(images=images, nsfw_content_detected=has_nsfw_concept)
@@ -18,17 +18,19 @@
import gc
import os
from collections import OrderedDict
from copy import copy
from typing import List, Optional, Union
from typing import List, Optional, Tuple, Union
import numpy as np
import onnx
import onnx_graphsurgeon as gs
import PIL.Image
import tensorrt as trt
import torch
from cuda import cudart
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from packaging import version
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
from polygraphy.backend.onnx.loader import fold_constants
@@ -40,23 +42,25 @@ from polygraphy.backend.trt import (
network_from_onnx_path,
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers import DiffusionPipeline
from diffusers.configuration_utils import FrozenDict, deprecate
from diffusers.image_processor import VaeImageProcessor
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
StableDiffusionPipeline,
StableDiffusionPipelineOutput,
StableDiffusionSafetyChecker,
)
from diffusers.schedulers import DDIMScheduler
from diffusers.utils import logging
from diffusers.utils.torch_utils import randn_tensor
"""
Installation instructions
python3 -m pip install --upgrade transformers diffusers>=0.16.0
python3 -m pip install --upgrade tensorrt>=8.6.1
python3 -m pip install --upgrade tensorrt~=10.2.0
python3 -m pip install --upgrade polygraphy>=0.47.0 onnx-graphsurgeon --extra-index-url https://pypi.ngc.nvidia.com
python3 -m pip install onnxruntime
"""
@@ -86,10 +90,6 @@ else:
torch_to_numpy_dtype_dict = {value: key for (key, value) in numpy_to_torch_dtype_dict.items()}
def device_view(t):
return cuda.DeviceView(ptr=t.data_ptr(), shape=t.shape, dtype=torch_to_numpy_dtype_dict[t.dtype])
class Engine:
def __init__(self, engine_path):
self.engine_path = engine_path
@@ -110,10 +110,8 @@ class Engine:
onnx_path,
fp16,
input_profile=None,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
workspace_size=0,
):
logger.warning(f"Building TensorRT engine for {onnx_path}: {self.engine_path}")
p = Profile()
@@ -122,20 +120,13 @@ class Engine:
assert len(dims) == 3
p.add(name, min=dims[0], opt=dims[1], max=dims[2])
config_kwargs = {}
config_kwargs["preview_features"] = [trt.PreviewFeature.DISABLE_EXTERNAL_TACTIC_SOURCES_FOR_CORE_0805]
if enable_preview:
# Faster dynamic shapes made optional since it increases engine build time.
config_kwargs["preview_features"].append(trt.PreviewFeature.FASTER_DYNAMIC_SHAPES_0805)
if workspace_size > 0:
config_kwargs["memory_pool_limits"] = {trt.MemoryPoolType.WORKSPACE: workspace_size}
extra_build_args = {}
if not enable_all_tactics:
config_kwargs["tactic_sources"] = []
extra_build_args["tactic_sources"] = []
engine = engine_from_network(
network_from_onnx_path(onnx_path, flags=[trt.OnnxParserFlag.NATIVE_INSTANCENORM]),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **config_kwargs),
config=CreateConfig(fp16=fp16, profiles=[p], load_timing_cache=timing_cache, **extra_build_args),
save_timing_cache=timing_cache,
)
save_engine(engine, path=self.engine_path)
@@ -148,28 +139,24 @@ class Engine:
self.context = self.engine.create_execution_context()
def allocate_buffers(self, shape_dict=None, device="cuda"):
for idx in range(trt_util.get_bindings_per_profile(self.engine)):
binding = self.engine[idx]
if shape_dict and binding in shape_dict:
shape = shape_dict[binding]
for binding in range(self.engine.num_io_tensors):
name = self.engine.get_tensor_name(binding)
if shape_dict and name in shape_dict:
shape = shape_dict[name]
else:
shape = self.engine.get_binding_shape(binding)
dtype = trt.nptype(self.engine.get_binding_dtype(binding))
if self.engine.binding_is_input(binding):
self.context.set_binding_shape(idx, shape)
shape = self.engine.get_tensor_shape(name)
dtype = trt.nptype(self.engine.get_tensor_dtype(name))
if self.engine.get_tensor_mode(name) == trt.TensorIOMode.INPUT:
self.context.set_input_shape(name, shape)
tensor = torch.empty(tuple(shape), dtype=numpy_to_torch_dtype_dict[dtype]).to(device=device)
self.tensors[binding] = tensor
self.buffers[binding] = cuda.DeviceView(ptr=tensor.data_ptr(), shape=shape, dtype=dtype)
self.tensors[name] = tensor
def infer(self, feed_dict, stream):
start_binding, end_binding = trt_util.get_active_profile_bindings(self.context)
# shallow copy of ordered dict
device_buffers = copy(self.buffers)
for name, buf in feed_dict.items():
assert isinstance(buf, cuda.DeviceView)
device_buffers[name] = buf
bindings = [0] * start_binding + [buf.ptr for buf in device_buffers.values()]
noerror = self.context.execute_async_v2(bindings=bindings, stream_handle=stream.ptr)
self.tensors[name].copy_(buf)
for name, tensor in self.tensors.items():
self.context.set_tensor_address(name, tensor.data_ptr())
noerror = self.context.execute_async_v3(stream)
if not noerror:
raise ValueError("ERROR: inference failed.")
@@ -310,10 +297,8 @@ def build_engines(
force_engine_rebuild=False,
static_batch=False,
static_shape=True,
enable_preview=False,
enable_all_tactics=False,
timing_cache=None,
max_workspace_size=0,
):
built_engines = {}
if not os.path.isdir(onnx_dir):
@@ -378,9 +363,7 @@ def build_engines(
static_batch=static_batch,
static_shape=static_shape,
),
enable_preview=enable_preview,
timing_cache=timing_cache,
workspace_size=max_workspace_size,
)
built_engines[model_name] = engine
@@ -588,11 +571,11 @@ def make_VAE(model, device, max_batch_size, embedding_dim, inpaint=False):
return VAE(model, device=device, max_batch_size=max_batch_size, embedding_dim=embedding_dim)
class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
class TensorRTStableDiffusionPipeline(DiffusionPipeline):
r"""
Pipeline for text-to-image generation using TensorRT accelerated Stable Diffusion.
This model inherits from [`StableDiffusionPipeline`]. Check the superclass documentation for the generic methods the
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
Args:
@@ -616,6 +599,8 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor"]
def __init__(
self,
vae: AutoencoderKL,
@@ -632,28 +617,90 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
image_width: int = 768,
max_batch_size: int = 16,
# ONNX export parameters
onnx_opset: int = 17,
onnx_opset: int = 18,
onnx_dir: str = "onnx",
# TensorRT engine build parameters
engine_dir: str = "engine",
build_preview_features: bool = True,
force_engine_rebuild: bool = False,
timing_cache: str = "timing_cache",
):
super().__init__(
vae,
text_encoder,
tokenizer,
unet,
scheduler,
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} is outdated. `steps_offset`"
f" should be set to 1 instead of {scheduler.config.steps_offset}. Please make sure "
"to update the config accordingly as leaving `steps_offset` might led to incorrect results"
" in future versions. If you have downloaded this checkpoint from the Hugging Face Hub,"
" it would be very nice if you could open a Pull request for the `scheduler/scheduler_config.json`"
" file"
)
deprecate("steps_offset!=1", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["steps_offset"] = 1
scheduler._internal_dict = FrozenDict(new_config)
if hasattr(scheduler.config, "clip_sample") and scheduler.config.clip_sample is True:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} has not set the configuration `clip_sample`."
" `clip_sample` should be set to False in the configuration file. Please make sure to update the"
" config accordingly as not setting `clip_sample` in the config might lead to incorrect results in"
" future versions. If you have downloaded this checkpoint from the Hugging Face Hub, it would be very"
" nice if you could open a Pull request for the `scheduler/scheduler_config.json` file"
)
deprecate("clip_sample not set", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["clip_sample"] = False
scheduler._internal_dict = FrozenDict(new_config)
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
version.parse(unet.config._diffusers_version).base_version
) < version.parse("0.9.0.dev0")
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
deprecation_message = (
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
" the `unet/config.json` file"
)
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(unet.config)
new_config["sample_size"] = 64
unet._internal_dict = FrozenDict(new_config)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
self.stages = stages
self.image_height, self.image_width = image_height, image_width
self.inpaint = False
@@ -664,7 +711,6 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
self.timing_cache = timing_cache
self.build_static_batch = False
self.build_dynamic_shape = False
self.build_preview_features = build_preview_features
self.max_batch_size = max_batch_size
# TODO: Restrict batch size to 4 for larger image dimensions as a WAR for TensorRT limitation.
@@ -675,6 +721,11 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
self.models = {} # loaded in __loadModels()
self.engine = {} # loaded in build_engines()
self.vae.forward = self.vae.decode
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor)
self.register_to_config(requires_safety_checker=requires_safety_checker)
def __loadModels(self):
# Load pipeline models
self.embedding_dim = self.text_encoder.config.hidden_size
@@ -691,6 +742,75 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
if "vae" in self.stages:
self.models["vae"] = make_VAE(self.vae, **models_args)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
def prepare_latents(
self,
batch_size: int,
num_channels_latents: int,
height: int,
width: int,
dtype: torch.dtype,
device: torch.device,
generator: Union[torch.Generator, List[torch.Generator]],
latents: Optional[torch.Tensor] = None,
) -> torch.Tensor:
r"""
Prepare the latent vectors for diffusion.
Args:
batch_size (int): The number of samples in the batch.
num_channels_latents (int): The number of channels in the latent vectors.
height (int): The height of the latent vectors.
width (int): The width of the latent vectors.
dtype (torch.dtype): The data type of the latent vectors.
device (torch.device): The device to place the latent vectors on.
generator (Union[torch.Generator, List[torch.Generator]]): The generator(s) to use for random number generation.
latents (Optional[torch.Tensor]): The pre-existing latent vectors. If None, new latent vectors will be generated.
Returns:
torch.Tensor: The prepared latent vectors.
"""
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
if latents is None:
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
else:
latents = latents.to(device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
return latents
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker
def run_safety_checker(
self, image: Union[torch.Tensor, PIL.Image.Image], device: torch.device, dtype: torch.dtype
) -> Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]:
r"""
Runs the safety checker on the given image.
Args:
image (Union[torch.Tensor, PIL.Image.Image]): The input image to be checked.
device (torch.device): The device to run the safety checker on.
dtype (torch.dtype): The data type of the input image.
Returns:
(image, has_nsfw_concept) Tuple[Union[torch.Tensor, PIL.Image.Image], Optional[bool]]: A tuple containing the processed image and
a boolean indicating whether the image has a NSFW (Not Safe for Work) concept.
"""
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
@@ -738,7 +858,6 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
force_engine_rebuild=self.force_engine_rebuild,
static_batch=self.build_static_batch,
static_shape=not self.build_dynamic_shape,
enable_preview=self.build_preview_features,
timing_cache=self.timing_cache,
)
@@ -769,9 +888,8 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
.to(self.torch_device)
)
text_input_ids_inp = device_view(text_input_ids)
# NOTE: output tensor for CLIP must be cloned because it will be overwritten when called again for negative prompt
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids_inp}, self.stream)[
text_embeddings = runEngine(self.engine["clip"], {"input_ids": text_input_ids}, self.stream)[
"text_embeddings"
].clone()
@@ -787,8 +905,7 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
.input_ids.type(torch.int32)
.to(self.torch_device)
)
uncond_input_ids_inp = device_view(uncond_input_ids)
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids_inp}, self.stream)[
uncond_embeddings = runEngine(self.engine["clip"], {"input_ids": uncond_input_ids}, self.stream)[
"text_embeddings"
]
@@ -812,18 +929,15 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
# Predict the noise residual
timestep_float = timestep.float() if timestep.dtype != torch.float32 else timestep
sample_inp = device_view(latent_model_input)
timestep_inp = device_view(timestep_float)
embeddings_inp = device_view(text_embeddings)
noise_pred = runEngine(
self.engine["unet"],
{"sample": sample_inp, "timestep": timestep_inp, "encoder_hidden_states": embeddings_inp},
{"sample": latent_model_input, "timestep": timestep_float, "encoder_hidden_states": text_embeddings},
self.stream,
)["latent"]
# Perform guidance
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + self.guidance_scale * (noise_pred_text - noise_pred_uncond)
noise_pred = noise_pred_uncond + self._guidance_scale * (noise_pred_text - noise_pred_uncond)
latents = self.scheduler.step(noise_pred, timestep, latents).prev_sample
@@ -831,12 +945,12 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
return latents
def __decode_latent(self, latents):
images = runEngine(self.engine["vae"], {"latent": device_view(latents)}, self.stream)["images"]
images = runEngine(self.engine["vae"], {"latent": latents}, self.stream)["images"]
images = (images / 2 + 0.5).clamp(0, 1)
return images.cpu().permute(0, 2, 3, 1).float().numpy()
def __loadResources(self, image_height, image_width, batch_size):
self.stream = cuda.Stream()
self.stream = cudart.cudaStreamCreate()[1]
# Allocate buffers for TensorRT engine bindings
for model_name, obj in self.models.items():
+2 -2
View File
@@ -425,8 +425,8 @@ pipe.load_lora_weights(lora_model_id)
image = pipe("A picture of a sks dog in a bucket", num_inference_steps=25).images[0]
```
Note that the use of [`LoraLoaderMixin.load_lora_weights`](https://huggingface.co/docs/diffusers/main/en/api/loaders#diffusers.loaders.LoraLoaderMixin.load_lora_weights) is preferred to [`UNet2DConditionLoadersMixin.load_attn_procs`](https://huggingface.co/docs/diffusers/main/en/api/loaders#diffusers.loaders.UNet2DConditionLoadersMixin.load_attn_procs) for loading LoRA parameters. This is because
`LoraLoaderMixin.load_lora_weights` can handle the following situations:
Note that the use of [`StableDiffusionLoraLoaderMixin.load_lora_weights`](https://huggingface.co/docs/diffusers/main/en/api/loaders#diffusers.loaders.StableDiffusionLoraLoaderMixin.load_lora_weights) is preferred to [`UNet2DConditionLoadersMixin.load_attn_procs`](https://huggingface.co/docs/diffusers/main/en/api/loaders#diffusers.loaders.UNet2DConditionLoadersMixin.load_attn_procs) for loading LoRA parameters. This is because
`StableDiffusionLoraLoaderMixin.load_lora_weights` can handle the following situations:
* LoRA parameters that don't have separate identifiers for the UNet and the text encoder (such as [`"patrickvonplaten/lora_dreambooth_dog_example"`](https://huggingface.co/patrickvonplaten/lora_dreambooth_dog_example)). So, you can just do:
+4 -4
View File
@@ -148,12 +148,12 @@ accelerate launch train_dreambooth_lora_sd3.py \
```
### Text Encoder Training
Alongside the transformer, LoRA fine-tuning of the CLIP text encoders is now also supported.
Alongside the transformer, LoRA fine-tuning of the CLIP text encoders is now also supported.
To do so, just specify `--train_text_encoder` while launching training. Please keep the following points in mind:
> [!NOTE]
> SD3 has three text encoders (CLIP L/14, OpenCLIP bigG/14, and T5-v1.1-XXL).
By enabling `--train_text_encoder`, LoRA fine-tuning of both **CLIP encoders** is performed. At the moment, T5 fine-tuning is not supported and weights remain frozen when text encoder training is enabled.
> SD3 has three text encoders (CLIP L/14, OpenCLIP bigG/14, and T5-v1.1-XXL).
By enabling `--train_text_encoder`, LoRA fine-tuning of both **CLIP encoders** is performed. At the moment, T5 fine-tuning is not supported and weights remain frozen when text encoder training is enabled.
To perform DreamBooth LoRA with text-encoder training, run:
```bash
@@ -185,4 +185,4 @@ accelerate launch train_dreambooth_lora_sd3.py \
1. We default to the "logit_normal" weighting scheme for the loss following the SD3 paper. Thanks to @bghira for helping us discover that for other weighting schemes supported from the training script, training may incur numerical instabilities.
2. Thanks to `bghira`, `JinxuXiang`, and `bendanzzc` for helping us discover a bug in how VAE encoding was being done previously. This has been fixed in [#8917](https://github.com/huggingface/diffusers/pull/8917).
3. Additionally, we now have the option to control if we want to apply preconditioning to the model outputs via a `--precondition_outputs` CLI arg. It affects how the model `target` is calculated as well.
3. Additionally, we now have the option to control if we want to apply preconditioning to the model outputs via a `--precondition_outputs` CLI arg. It affects how the model `target` is calculated as well.
+4 -4
View File
@@ -52,7 +52,7 @@ from diffusers import (
StableDiffusionPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params
from diffusers.utils import (
@@ -956,7 +956,7 @@ def main(args):
# make sure to pop weight so that corresponding model is not saved again
weights.pop()
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
output_dir,
unet_lora_layers=unet_lora_layers_to_save,
text_encoder_lora_layers=text_encoder_lora_layers_to_save,
@@ -976,7 +976,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
@@ -1376,7 +1376,7 @@ def main(args):
else:
text_encoder_state_dict = None
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_state_dict,
text_encoder_lora_layers=text_encoder_state_dict,
@@ -58,7 +58,7 @@ from diffusers import (
StableDiffusionXLPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params, compute_snr
from diffusers.utils import (
@@ -1260,7 +1260,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
@@ -49,7 +49,7 @@ from diffusers import (
DPMSolverMultistepScheduler,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.utils import check_min_version, convert_state_dict_to_diffusers
from diffusers.utils.import_utils import is_xformers_available
@@ -604,7 +604,7 @@ def main(args):
# make sure to pop weight so that corresponding model is not saved again
weights.pop()
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
output_dir,
unet_lora_layers=unet_lora_layers_to_save,
text_encoder_lora_layers=None,
@@ -621,8 +621,8 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
LoraLoaderMixin.load_lora_into_unet(lora_state_dict, network_alphas=network_alphas, unet=unet_)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
StableDiffusionLoraLoaderMixin.load_lora_into_unet(lora_state_dict, network_alphas=network_alphas, unet=unet_)
accelerator.register_save_state_pre_hook(save_model_hook)
accelerator.register_load_state_pre_hook(load_model_hook)
@@ -951,7 +951,7 @@ def main(args):
unet = unet.to(torch.float32)
unet_lora_state_dict = convert_state_dict_to_diffusers(get_peft_model_state_dict(unet))
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
save_directory=args.output_dir, unet_lora_layers=unet_lora_state_dict, text_encoder_lora_layers=None
)
@@ -46,5 +46,4 @@ pipe.enable_model_cpu_offload()
# generate image
generator = torch.manual_seed(0)
image = pipe("a tortoise", num_inference_steps=20, generator=generator, image_pair=[image_a,image_b], image=query).images[0]
```
@@ -2051,7 +2051,7 @@ if __name__ == "__main__":
default=512,
type=int,
help=(
"The image size that the model was trained on. Use 512 for Stable Diffusion v1.X and Stable Siffusion v2"
"The image size that the model was trained on. Use 512 for Stable Diffusion v1.X and Stable Diffusion v2"
" Base. Use 768 for Stable Diffusion v2."
),
)
@@ -28,7 +28,7 @@ import torch.nn.functional as F
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.loaders import FromSingleFileMixin, StableDiffusionLoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, ControlNetModel, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.controlnet.multicontrolnet import MultiControlNetModel
@@ -142,7 +142,9 @@ def retrieve_timesteps(
return timesteps, num_inference_steps
class PromptDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin):
class PromptDiffusionPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, StableDiffusionLoraLoaderMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion with ControlNet guidance.
@@ -153,8 +155,8 @@ class PromptDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.StableDiffusionLoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
@@ -348,7 +350,7 @@ class PromptDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
if lora_scale is not None and isinstance(self, StableDiffusionLoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
@@ -480,7 +482,7 @@ class PromptDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
if isinstance(self, StableDiffusionLoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
@@ -1251,7 +1253,7 @@ class PromptDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
)
if guess_mode and self.do_classifier_free_guidance:
# Infered ControlNet only for the conditional batch.
# Inferred ControlNet only for the conditional batch.
# To apply the output of ControlNet to both the unconditional and conditional batches,
# add 0 to the unconditional batch to keep it unchanged.
down_block_res_samples = [torch.cat([torch.zeros_like(d), d]) for d in down_block_res_samples]
@@ -52,7 +52,7 @@ from diffusers import (
StableDiffusionPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params
from diffusers.utils import (
@@ -999,7 +999,7 @@ def main(args):
# make sure to pop weight so that corresponding model is not saved again
weights.pop()
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
output_dir,
unet_lora_layers=unet_lora_layers_to_save,
text_encoder_lora_layers=text_encoder_lora_layers_to_save,
@@ -1019,7 +1019,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
@@ -1451,7 +1451,7 @@ def main(args):
else:
text_encoder_state_dict = None
LoraLoaderMixin.save_lora_weights(
StableDiffusionLoraLoaderMixin.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_state_dict,
text_encoder_lora_layers=text_encoder_state_dict,
@@ -59,7 +59,7 @@ from diffusers import (
StableDiffusionXLPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params, compute_snr
from diffusers.utils import (
@@ -1334,7 +1334,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, network_alphas = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, network_alphas = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
@@ -49,7 +49,7 @@ from diffusers import (
StableDiffusionXLPipeline,
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.loaders import StableDiffusionLoraLoaderMixin
from diffusers.optimization import get_scheduler
from diffusers.training_utils import _set_state_dict_into_text_encoder, cast_training_params, compute_snr
from diffusers.utils import (
@@ -749,7 +749,7 @@ def main(args):
else:
raise ValueError(f"unexpected save model: {model.__class__}")
lora_state_dict, _ = LoraLoaderMixin.lora_state_dict(input_dir)
lora_state_dict, _ = StableDiffusionLoraLoaderMixin.lora_state_dict(input_dir)
unet_state_dict = {f'{k.replace("unet.", "")}': v for k, v in lora_state_dict.items() if k.startswith("unet.")}
unet_state_dict = convert_unet_state_dict_to_peft(unet_state_dict)
incompatible_keys = set_peft_model_state_dict(unet_, unet_state_dict, adapter_name="default")
@@ -11,28 +11,28 @@ huggingface-cli login
This will also allow us to push the trained model parameters to the Hugging Face Hub platform.
For setup, inference code, and details on how to run the code, please follow the Colab Notebook provided above.
For setup, inference code, and details on how to run the code, please follow the Colab Notebook provided above.
## How
We make use of several techniques to make this possible:
* Compute the embeddings from the instance prompt and serialize them for later reuse. This is implemented in the [`compute_embeddings.py`](./compute_embeddings.py) script. We use an 8bit (as introduced in [`LLM.int8()`](https://arxiv.org/abs/2208.07339)) T5 to reduce memory requirements to ~10.5GB.
* Compute the embeddings from the instance prompt and serialize them for later reuse. This is implemented in the [`compute_embeddings.py`](./compute_embeddings.py) script. We use an 8bit (as introduced in [`LLM.int8()`](https://arxiv.org/abs/2208.07339)) T5 to reduce memory requirements to ~10.5GB.
* In the `train_dreambooth_sd3_lora_miniature.py` script, we make use of:
* 8bit Adam for optimization through the `bitsandbytes` library.
* Gradient checkpointing and gradient accumulation.
* FP16 precision.
* Flash attention through `F.scaled_dot_product_attention()`.
* Flash attention through `F.scaled_dot_product_attention()`.
Computing the text embeddings is arguably the most memory-intensive part in the pipeline as SD3 employs three text encoders. If we run them in FP32, it will take about 20GB of VRAM. With FP16, we are down to 12GB.
Computing the text embeddings is arguably the most memory-intensive part in the pipeline as SD3 employs three text encoders. If we run them in FP32, it will take about 20GB of VRAM. With FP16, we are down to 12GB.
## Gotchas
This project is educational. It exists to showcase the possibility of fine-tuning a big diffusion system on consumer GPUs. But additional components might have to be added to obtain state-of-the-art performance. Below are some commonly known gotchas that users should be aware of:
* Training of text encoders is purposefully disabled.
* Techniques such as prior-preservation is unsupported.
* Training of text encoders is purposefully disabled.
* Techniques such as prior-preservation is unsupported.
* Custom instance captions for instance images are unsupported, but this should be relatively easy to integrate.
Hopefully, this project gives you a template to extend it further to suit your needs.

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