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23 Commits

Author SHA1 Message Date
patil-suraj ea238e821b up 2024-03-18 11:47:47 +01:00
patil-suraj b6d1d670fc up 2024-03-18 11:34:17 +01:00
Dhruv Nair 4330a747d4 [Tests] Fix ControlNet Single File tests (#7315)
* update

* update

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-18 11:28:59 +05:30
Sayak Paul 76de6a09fb post-release v0.27.0 (#7329)
* post-release

* quality
2024-03-18 10:52:20 +05:30
Sayak Paul 25caf24ef9 Fix release workflow deps (#7339)
* pop scale from the top-level unet instead of getting it.

* improve readability.

* fix: pypi workflow deps

* revert
2024-03-16 07:18:11 +05:30
Abubakar Abid 8db3c9bc9f Adds docs for gradio.Interface.from_pipeline() (#7346)
* gradio docs

* Update docs/source/en/api/pipelines/stable_diffusion/overview.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* changes

* changes

* changes

* Update docs/source/en/api/pipelines/stable_diffusion/overview.md

---------

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-16 07:11:28 +05:30
Sayak Paul e0e9f81971 add: torch to the pypi step. (#7328) 2024-03-15 12:28:12 +05:30
M. Tolga Cangöz 5d848ec07c [Tests] Update a deprecated parameter in test files and fix several typos (#7277)
* Add properties and `IPAdapterTesterMixin` tests for `StableDiffusionPanoramaPipeline`

* Fix variable name typo and update comments

* Update deprecated `output_type="numpy"` to "np" in test files

* Discard changes to src/diffusers/pipelines/stable_diffusion_panorama/pipeline_stable_diffusion_panorama.py

* Update test_stable_diffusion_panorama.py

* Update numbers in README.md

* Update get_guidance_scale_embedding method to use timesteps instead of w

* Update number of checkpoints in README.md

* Add type hints and fix var name

* Fix PyTorch's convention for inplace functions

* Fix a typo

* Revert "Fix PyTorch's convention for inplace functions"

This reverts commit 74350cf65b.

* Fix typos

* Indent

* Refactor get_guidance_scale_embedding method in LEditsPPPipelineStableDiffusionXL class
2024-03-14 12:17:35 -07:00
Dhruv Nair 4974b84564 Update Cascade Tests (#7324)
* update

* update

* update
2024-03-14 20:51:22 +05:30
Linoy Tsaban 83062fb872 [Advanced DreamBooth LoRA SDXL] Support EDM-style training (follow up of #7126) (#7182)
* add edm style training

* style

* finish adding edm training feature

* import fix

* fix latents mean

* minor adjustments

* add edm to readme

* style

* fix autocast and scheduler config issues when using edm

* style

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-14 18:40:14 +05:30
Suraj Patil b6d7e31d10 add edm schedulers in doc (#7319)
* add edm schedulers in doc

* add in toctree

* address reviewe comments
2024-03-14 11:52:25 +01:00
Anatoly Belikov 53e9aacc10 log loss per image (#7278)
* log loss per image

* add commandline param for per image loss logging

* style

* debug-loss -> debug_loss

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-14 11:41:43 +05:30
Dhruv Nair 41424466e3 [Tests] Fix incorrect constant in VAE scaling test. (#7301)
update
2024-03-14 10:24:01 +05:30
Sayak Paul 95de1981c9 add: pytest log installation (#7313) 2024-03-14 10:01:16 +05:30
Kenneth Gerald Hamilton 0b45b58867 update get_order_list if statement (#7309)
* update get_order_list if statement

* revery
2024-03-13 18:29:42 -10:00
Beinsezii d3986f18be Change step_offset scheduler docstrings (#7128)
* Change step_offset scheduler docstrings

* Mention it may be needed by some models

* More docstrings

These ones failed literal S&R because I performed it case-sensitive
which is fun.

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-13 15:12:00 -10:00
Alexander Bonnet ee6a3a993d Fix typos in UNet2DConditionModel documentation (#7291)
* fix typo in UNet2DConditionModel documentation

* Fix indentation that may fix doc rendering

* Fix squished doc lines
2024-03-13 09:31:29 -07:00
Michael b300517305 Add Intro page of TCD (#7259)
* add tcd intro

* resolve repos

* Apply suggestions from code review

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* revise NFEs related

* change inpainting location

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2024-03-13 09:21:51 -07:00
jnhuang ac07b6dc6a Fix Wrong Text-encoder Grad Setting in Custom_Diffusion Training (#7302)
fix index in set textencoder grad

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-13 20:22:44 +05:30
Sayak Paul 46ab56a468 add: support for notifying maintainers about the nightly test status (#7117)
* add: support for notifying maintainers about the nightly test status

* add: a tempoerary workflow for validation.

* cancel in progress.

* runs-on

* clean up

* add: peft dep

* change device.

* multiple edits.

* remove temp workflow.
2024-03-13 16:48:11 +05:30
Sayak Paul 038ff70023 [PyPI publishing] feat: automate the process of pypi publication to some extent. (#7270)
* feat: automate the process of pypi publication to some extent.

* utility to fetch the latest release branch

* correct package name.
2024-03-13 16:27:59 +05:30
Manuel Brack 00eca4b887 [Pipeline] Add LEDITS++ pipelines (#6074)
* Setup LEdits++ file structure

* Fix import

* LEditsPP Stable Diffusion pipeline

* Include variable image aspect ratios

* Implement LEDITS++ for SDXL

* clean up LEditsPPPipelineStableDiffusion

* Adjust inversion output

* Added docu, more cleanup for LEditsPPPipelineStableDiffusion

* clean up LEditsPPPipelineStableDiffusionXL

* Update documentation

* Fix documentation import

* Add skeleton IF implementation

* Fix documentation typo

* Add LEDTIS docu to toctree

* Add missing title

* Finalize SD documentation

* Finalize SD-XL documentation

* Fix code style and quality

* Fix typo

* Fix return types

* added LEditsPPPipelineIF; minor changes for LEditsPPPipelineStableDiffusion and LEditsPPPipelineStableDiffusionXL

* Fix copy reference

* add documentation for IF

* Add first tests

* Fix batching for SD-XL

* Fix text encoding and perfect reconstruction for SD-XL

* Add tests for SD-XL, minor changes

* move user_mask to correct device, use cross_attention_kwargs also for inversion

* Example docstring

* Fix attention resolution for non-square images

* Refactoring for PR review

* Safely remove ledits_utils.py

* Style fixes

* Replace assertions with ValueError

* Remove LEditsPPPipelineIF

* Remove unecessary input checks

* Refactoring of CrossAttnProcessor

* Revert unecessary changes to scheduler

* Remove first progress-bar in inversion

* Refactor scheduler usage and reset

* Use imageprocessor instead of custom logic

* Fix scheduler init warning

* Fix error when running the pipeline in fp16

* Update documentation wrt perfect inversion

* Update tests

* Fix code quality and copy consistency

* Update LEditsPP import

* Remove enable/disable methods that are now in StableDiffusionMixin

* Change import in docs

* Revert import structure change

* Fix ledits imports

---------

Co-authored-by: Katharina Kornmeier <katharina.kornmeier@stud.tu-darmstadt.de>
2024-03-13 12:43:47 +02:00
Dhruv Nair 30132aba30 Update Stable Cascade Conversion Scripts (#7271)
* update

* update

* update

* update

* update

* update

* update

* update

* update

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-03-13 12:35:44 +05:30
160 changed files with 5617 additions and 430 deletions
+22 -2
View File
@@ -12,6 +12,7 @@ env:
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
RUN_NIGHTLY: yes
SLACK_API_TOKEN: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
jobs:
run_nightly_tests:
@@ -64,6 +65,7 @@ jobs:
python -m uv pip install -e [quality,test]
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
python -m uv pip install pytest-reportlog
- name: Environment
run: |
@@ -78,7 +80,8 @@ jobs:
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
--report-log=${{ matrix.config.report }}.log \
tests/
- name: Run nightly Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
@@ -89,6 +92,7 @@ jobs:
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
--report-log=${{ matrix.config.report }}.log \
tests/
- name: Run nightly ONNXRuntime CUDA tests
@@ -100,6 +104,7 @@ jobs:
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
--report-log=${{ matrix.config.report }}.log \
tests/
- name: Failure short reports
@@ -112,6 +117,12 @@ jobs:
with:
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
@@ -140,6 +151,7 @@ jobs:
${CONDA_RUN} python -m uv pip install -e [quality,test]
${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m uv pip install pytest-reportlog
- name: Environment
shell: arch -arch arm64 bash {0}
@@ -152,7 +164,9 @@ jobs:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
--report-log=tests_torch_mps.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
@@ -164,3 +178,9 @@ jobs:
with:
name: torch_mps_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
@@ -0,0 +1,23 @@
name: Notify Slack about a release
on:
workflow_dispatch:
release:
types: [published]
jobs:
build:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Notify Slack about the release
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL }}
run: pip install requests && python utils/notify_slack_about_release.py
+81
View File
@@ -0,0 +1,81 @@
# Adapted from https://blog.deepjyoti30.dev/pypi-release-github-action
name: PyPI release
on:
workflow_dispatch:
push:
tags:
- "*"
jobs:
find-and-checkout-latest-branch:
runs-on: ubuntu-latest
outputs:
latest_branch: ${{ steps.set_latest_branch.outputs.latest_branch }}
steps:
- name: Checkout Repo
uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Fetch latest branch
id: fetch_latest_branch
run: |
pip install -U requests packaging
LATEST_BRANCH=$(python utils/fetch_latest_release_branch.py)
echo "Latest branch: $LATEST_BRANCH"
echo "latest_branch=$LATEST_BRANCH" >> $GITHUB_ENV
- name: Set latest branch output
id: set_latest_branch
run: echo "::set-output name=latest_branch::${{ env.latest_branch }}"
release:
needs: find-and-checkout-latest-branch
runs-on: ubuntu-latest
steps:
- name: Checkout Repo
uses: actions/checkout@v3
with:
ref: ${{ needs.find-and-checkout-latest-branch.outputs.latest_branch }}
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -U setuptools wheel twine
pip install -U torch --index-url https://download.pytorch.org/whl/cpu
pip install -U transformers
- name: Build the dist files
run: python setup.py bdist_wheel && python setup.py sdist
- name: Publish to the test PyPI
env:
TWINE_USERNAME: ${{ secrets.TEST_PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.TEST_PYPI_PASSWORD }}
run: twine upload dist/* -r pypitest --repository-url=https://test.pypi.org/legacy/
- name: Test installing diffusers and importing
run: |
pip install diffusers && pip uninstall diffusers -y
pip install -i https://testpypi.python.org/pypi diffusers
python -c "from diffusers import __version__; print(__version__)"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('fusing/unet-ldm-dummy-update'); pipe()"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('hf-internal-testing/tiny-stable-diffusion-pipe', safety_checker=None); pipe('ah suh du')"
python -c "from diffusers import *"
- name: Publish to PyPI
env:
TWINE_USERNAME: ${{ secrets.PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.PYPI_PASSWORD }}
run: twine upload dist/* -r pypi
+2 -2
View File
@@ -77,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 19000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 22000+ checkpoints):
```python
from diffusers import DiffusionPipeline
@@ -219,7 +219,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +8000 other amazing GitHub repositories 💪
- +9000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
+8
View File
@@ -104,6 +104,8 @@
title: Latent Consistency Model-LoRA
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/inference_with_tcd_lora
title: Trajectory Consistency Distillation-LoRA
- local: using-diffusers/svd
title: Stable Video Diffusion
title: Specific pipeline examples
@@ -304,6 +306,8 @@
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
@@ -396,6 +400,10 @@
title: DPMSolverSDEScheduler
- local: api/schedulers/singlestep_dpm_solver
title: DPMSolverSinglestepScheduler
- local: api/schedulers/edm_multistep_dpm_solver
title: EDMDPMSolverMultistepScheduler
- local: api/schedulers/edm_euler
title: EDMEulerScheduler
- local: api/schedulers/euler_ancestral
title: EulerAncestralDiscreteScheduler
- local: api/schedulers/euler
+54
View File
@@ -0,0 +1,54 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# LEDITS++
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:
*Text-to-image diffusion models have recently received increasing interest for their astonishing ability to produce high-fidelity images from solely text inputs. Subsequent research efforts aim to exploit and apply their capabilities to real image editing. However, existing image-to-image methods are often inefficient, imprecise, and of limited versatility. They either require time-consuming fine-tuning, deviate unnecessarily strongly from the input image, and/or lack support for multiple, simultaneous edits. To address these issues, we introduce LEDITS++, an efficient yet versatile and precise textual image manipulation technique. LEDITS++'s novel inversion approach requires no tuning nor optimization and produces high-fidelity results with a few diffusion steps. Second, our methodology supports multiple simultaneous edits and is architecture-agnostic. Third, we use a novel implicit masking technique that limits changes to relevant image regions. We propose the novel TEdBench++ benchmark as part of our exhaustive evaluation. Our results demonstrate the capabilities of LEDITS++ and its improvements over previous methods. The project page is available at https://leditsplusplus-project.static.hf.space .*
<Tip>
You can find additional information about LEDITS++ on the [project page](https://leditsplusplus-project.static.hf.space/index.html) and try it out in a [demo](https://huggingface.co/spaces/editing-images/leditsplusplus).
</Tip>
<Tip warning={true}>
Due to some backward compatability issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
This issue is unlikely to have any noticeable effects on applied use-cases. However, we provide an alternative implementation that guarantees perfect inversion in a dedicated [GitHub repo](https://github.com/ml-research/ledits_pp).
</Tip>
We provide two distinct pipelines based on different pre-trained models.
## LEditsPPPipelineStableDiffusion
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusion
- all
- __call__
- invert
## LEditsPPPipelineStableDiffusionXL
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusionXL
- all
- __call__
- invert
## LEditsPPDiffusionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPDiffusionPipelineOutput
- all
## LEditsPPInversionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPInversionPipelineOutput
- all
+1
View File
@@ -57,6 +57,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latent Consistency Models](latent_consistency_models) | text2image |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D, text-to-pano, upscaling |
| [LEDITS++](ledits_pp) | image editing |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [Paint by Example](paint_by_example) | inpainting |
@@ -30,6 +30,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## StableDiffusionSafePipelineOutput
## SemanticStableDiffusionPipelineOutput
[[autodoc]] pipelines.semantic_stable_diffusion.pipeline_output.SemanticStableDiffusionPipelineOutput
- all
@@ -172,3 +172,41 @@ inpaint = StableDiffusionInpaintPipeline(**text2img.components)
# now you can use text2img(...), img2img(...), inpaint(...) just like the call methods of each respective pipeline
```
### Create web demos using `gradio`
The Stable Diffusion pipelines are automatically supported in [Gradio](https://github.com/gradio-app/gradio/), a library that makes creating beautiful and user-friendly machine learning apps on the web a breeze. First, make sure you have Gradio installed:
```
pip install -U gradio
```
Then, create a web demo around any Stable Diffusion-based pipeline. For example, you can create an image generation pipeline in a single line of code with Gradio's [`Interface.from_pipeline`](https://www.gradio.app/docs/interface#interface-from-pipeline) function:
```py
from diffusers import StableDiffusionPipeline
import gradio as gr
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
gr.Interface.from_pipeline(pipe).launch()
```
which opens an intuitive drag-and-drop interface in your browser:
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/gradio-panda.png)
Similarly, you could create a demo for an image-to-image pipeline with:
```py
from diffusers import StableDiffusionImg2ImgPipeline
import gradio as gr
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
gr.Interface.from_pipeline(pipe).launch()
```
By default, the web demo runs on a local server. If you'd like to share it with others, you can generate a temporary public
link by setting `share=True` in `launch()`. Or, you can host your demo on [Hugging Face Spaces](https://huggingface.co/spaces)https://huggingface.co/spaces for a permanent link.
@@ -0,0 +1,22 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMEulerScheduler
The Karras formulation of the Euler scheduler (Algorithm 2) from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper by Karras et al. This is a fast scheduler which can often generate good outputs in 20-30 steps. The scheduler is based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by [Katherine Crowson](https://github.com/crowsonkb/).
## EDMEulerScheduler
[[autodoc]] EDMEulerScheduler
## EDMEulerSchedulerOutput
[[autodoc]] schedulers.scheduling_edm_euler.EDMEulerSchedulerOutput
@@ -0,0 +1,24 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMDPMSolverMultistepScheduler
`EDMDPMSolverMultistepScheduler` is a [Karras formulation](https://huggingface.co/papers/2206.00364) of `DPMSolverMultistep`, a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
DPMSolver (and the improved version DPMSolver++) is a fast dedicated high-order solver for diffusion ODEs with convergence order guarantee. Empirically, DPMSolver sampling with only 20 steps can generate high-quality
samples, and it can generate quite good samples even in 10 steps.
## EDMDPMSolverMultistepScheduler
[[autodoc]] EDMDPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
@@ -0,0 +1,438 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
[[open-in-colab]]
# Trajectory Consistency Distillation-LoRA
Trajectory Consistency Distillation (TCD) enables a model to generate higher quality and more detailed images with fewer steps. Moreover, owing to the effective error mitigation during the distillation process, TCD demonstrates superior performance even under conditions of large inference steps.
The major advantages of TCD are:
- Better than Teacher: TCD demonstrates superior generative quality at both small and large inference steps and exceeds the performance of [DPM-Solver++(2S)](../../api/schedulers/multistep_dpm_solver) with Stable Diffusion XL (SDXL). There is no additional discriminator or LPIPS supervision included during TCD training.
- Flexible Inference Steps: The inference steps for TCD sampling can be freely adjusted without adversely affecting the image quality.
- Freely change detail level: During inference, the level of detail in the image can be adjusted with a single hyperparameter, *gamma*.
> [!TIP]
> For more technical details of TCD, please refer to the [paper](https://arxiv.org/abs/2402.19159) or official [project page](https://mhh0318.github.io/tcd/)).
For large models like SDXL, TCD is trained with [LoRA](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) to reduce memory usage. This is also useful because you can reuse LoRAs between different finetuned models, as long as they share the same base model, without further training.
This guide will show you how to perform inference with TCD-LoRAs for a variety of tasks like text-to-image and inpainting, as well as how you can easily combine TCD-LoRAs with other adapters. Choose one of the supported base model and it's corresponding TCD-LoRA checkpoint from the table below to get started.
| Base model | TCD-LoRA checkpoint |
|-------------------------------------------------------------------------------------------------|----------------------------------------------------------------|
| [stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) | [TCD-SD15](https://huggingface.co/h1t/TCD-SD15-LoRA) |
| [stable-diffusion-2-1-base](https://huggingface.co/stabilityai/stable-diffusion-2-1-base) | [TCD-SD21-base](https://huggingface.co/h1t/TCD-SD21-base-LoRA) |
| [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | [TCD-SDXL](https://huggingface.co/h1t/TCD-SDXL-LoRA) |
Make sure you have [PEFT](https://github.com/huggingface/peft) installed for better LoRA support.
```bash
pip install -U peft
```
## General tasks
In this guide, let's use the [`StableDiffusionXLPipeline`] and the [`TCDScheduler`]. Use the [`~StableDiffusionPipeline.load_lora_weights`] method to load the SDXL-compatible TCD-LoRA weights.
A few tips to keep in mind for TCD-LoRA inference are to:
- Keep the `num_inference_steps` between 4 and 50
- Set `eta` (used to control stochasticity at each step) between 0 and 1. You should use a higher `eta` when increasing the number of inference steps, but the downside is that a larger `eta` in [`TCDScheduler`] leads to blurrier images. A value of 0.3 is recommended to produce good results.
<hfoptions id="tasks">
<hfoption id="text-to-image">
```python
import torch
from diffusers import StableDiffusionXLPipeline, TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "Painting of the orange cat Otto von Garfield, Count of Bismarck-Schönhausen, Duke of Lauenburg, Minister-President of Prussia. Depicted wearing a Prussian Pickelhaube and eating his favorite meal - lasagna."
image = pipe(
prompt=prompt,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/demo_image.png)
</hfoption>
<hfoption id="inpainting">
```python
import torch
from diffusers import AutoPipelineForInpainting, TCDScheduler
from diffusers.utils import load_image, make_image_grid
device = "cuda"
base_model_id = "diffusers/stable-diffusion-xl-1.0-inpainting-0.1"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = AutoPipelineForInpainting.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = load_image(img_url).resize((1024, 1024))
mask_image = load_image(mask_url).resize((1024, 1024))
prompt = "a tiger sitting on a park bench"
image = pipe(
prompt=prompt,
image=init_image,
mask_image=mask_image,
num_inference_steps=8,
guidance_scale=0,
eta=0.3,
strength=0.99, # make sure to use `strength` below 1.0
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/inpainting_tcd.png)
</hfoption>
</hfoptions>
## Community models
TCD-LoRA also works with many community finetuned models and plugins. For example, load the [animagine-xl-3.0](https://huggingface.co/cagliostrolab/animagine-xl-3.0) checkpoint which is a community finetuned version of SDXL for generating anime images.
```python
import torch
from diffusers import StableDiffusionXLPipeline, TCDScheduler
device = "cuda"
base_model_id = "cagliostrolab/animagine-xl-3.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "A man, clad in a meticulously tailored military uniform, stands with unwavering resolve. The uniform boasts intricate details, and his eyes gleam with determination. Strands of vibrant, windswept hair peek out from beneath the brim of his cap."
image = pipe(
prompt=prompt,
num_inference_steps=8,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/animagine_xl.png)
TCD-LoRA also supports other LoRAs trained on different styles. For example, let's load the [TheLastBen/Papercut_SDXL](https://huggingface.co/TheLastBen/Papercut_SDXL) LoRA and fuse it with the TCD-LoRA with the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method.
> [!TIP]
> Check out the [Merge LoRAs](merge_loras) guide to learn more about efficient merging methods.
```python
import torch
from diffusers import StableDiffusionXLPipeline
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
styled_lora_id = "TheLastBen/Papercut_SDXL"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id, adapter_name="tcd")
pipe.load_lora_weights(styled_lora_id, adapter_name="style")
pipe.set_adapters(["tcd", "style"], adapter_weights=[1.0, 1.0])
prompt = "papercut of a winter mountain, snow"
image = pipe(
prompt=prompt,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/styled_lora.png)
## Adapters
TCD-LoRA is very versatile, and it can be combined with other adapter types like ControlNets, IP-Adapter, and AnimateDiff.
<hfoptions id="adapters">
<hfoption id="ControlNet">
### Depth ControlNet
```python
import torch
import numpy as np
from PIL import Image
from transformers import DPTFeatureExtractor, DPTForDepthEstimation
from diffusers import ControlNetModel, StableDiffusionXLControlNetPipeline
from diffusers.utils import load_image, make_image_grid
from scheduling_tcd import TCDScheduler
device = "cuda"
depth_estimator = DPTForDepthEstimation.from_pretrained("Intel/dpt-hybrid-midas").to(device)
feature_extractor = DPTFeatureExtractor.from_pretrained("Intel/dpt-hybrid-midas")
def get_depth_map(image):
image = feature_extractor(images=image, return_tensors="pt").pixel_values.to(device)
with torch.no_grad(), torch.autocast(device):
depth_map = depth_estimator(image).predicted_depth
depth_map = torch.nn.functional.interpolate(
depth_map.unsqueeze(1),
size=(1024, 1024),
mode="bicubic",
align_corners=False,
)
depth_min = torch.amin(depth_map, dim=[1, 2, 3], keepdim=True)
depth_max = torch.amax(depth_map, dim=[1, 2, 3], keepdim=True)
depth_map = (depth_map - depth_min) / (depth_max - depth_min)
image = torch.cat([depth_map] * 3, dim=1)
image = image.permute(0, 2, 3, 1).cpu().numpy()[0]
image = Image.fromarray((image * 255.0).clip(0, 255).astype(np.uint8))
return image
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
controlnet_id = "diffusers/controlnet-depth-sdxl-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
controlnet = ControlNetModel.from_pretrained(
controlnet_id,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_id,
controlnet=controlnet,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe.enable_model_cpu_offload()
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "stormtrooper lecture, photorealistic"
image = load_image("https://huggingface.co/lllyasviel/sd-controlnet-depth/resolve/main/images/stormtrooper.png")
depth_image = get_depth_map(image)
controlnet_conditioning_scale = 0.5 # recommended for good generalization
image = pipe(
prompt,
image=depth_image,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
controlnet_conditioning_scale=controlnet_conditioning_scale,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([depth_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/controlnet_depth_tcd.png)
### Canny ControlNet
```python
import torch
from diffusers import ControlNetModel, StableDiffusionXLControlNetPipeline
from diffusers.utils import load_image, make_image_grid
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
controlnet_id = "diffusers/controlnet-canny-sdxl-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
controlnet = ControlNetModel.from_pretrained(
controlnet_id,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_id,
controlnet=controlnet,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe.enable_model_cpu_offload()
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "ultrarealistic shot of a furry blue bird"
canny_image = load_image("https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny.png")
controlnet_conditioning_scale = 0.5 # recommended for good generalization
image = pipe(
prompt,
image=canny_image,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
controlnet_conditioning_scale=controlnet_conditioning_scale,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([canny_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/controlnet_canny_tcd.png)
<Tip>
The inference parameters in this example might not work for all examples, so we recommend you to try different values for `num_inference_steps`, `guidance_scale`, `controlnet_conditioning_scale` and `cross_attention_kwargs` parameters and choose the best one.
</Tip>
</hfoption>
<hfoption id="IP-Adapter">
This example shows how to use the TCD-LoRA with the [IP-Adapter](https://github.com/tencent-ailab/IP-Adapter/tree/main) and SDXL.
```python
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers.utils import load_image, make_image_grid
from ip_adapter import IPAdapterXL
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_path = "stabilityai/stable-diffusion-xl-base-1.0"
image_encoder_path = "sdxl_models/image_encoder"
ip_ckpt = "sdxl_models/ip-adapter_sdxl.bin"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(
base_model_path,
torch_dtype=torch.float16,
variant="fp16"
)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
ip_model = IPAdapterXL(pipe, image_encoder_path, ip_ckpt, device)
ref_image = load_image("https://raw.githubusercontent.com/tencent-ailab/IP-Adapter/main/assets/images/woman.png").resize((512, 512))
prompt = "best quality, high quality, wearing sunglasses"
image = ip_model.generate(
pil_image=ref_image,
prompt=prompt,
scale=0.5,
num_samples=1,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
seed=0,
)[0]
grid_image = make_image_grid([ref_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/ip_adapter.png)
</hfoption>
<hfoption id="AnimateDiff">
[`AnimateDiff`] allows animating images using Stable Diffusion models. TCD-LoRA can substantially accelerate the process without degrading image quality. The quality of animation with TCD-LoRA and AnimateDiff has a more lucid outcome.
```python
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from scheduling_tcd import TCDScheduler
from diffusers.utils import export_to_gif
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5")
pipe = AnimateDiffPipeline.from_pretrained(
"frankjoshua/toonyou_beta6",
motion_adapter=adapter,
).to("cuda")
# set TCDScheduler
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
# load TCD LoRA
pipe.load_lora_weights("h1t/TCD-SD15-LoRA", adapter_name="tcd")
pipe.load_lora_weights("guoyww/animatediff-motion-lora-zoom-in", weight_name="diffusion_pytorch_model.safetensors", adapter_name="motion-lora")
pipe.set_adapters(["tcd", "motion-lora"], adapter_weights=[1.0, 1.2])
prompt = "best quality, masterpiece, 1girl, looking at viewer, blurry background, upper body, contemporary, dress"
generator = torch.manual_seed(0)
frames = pipe(
prompt=prompt,
num_inference_steps=5,
guidance_scale=0,
cross_attention_kwargs={"scale": 1},
num_frames=24,
eta=0.3,
generator=generator
).frames[0]
export_to_gif(frames, "animation.gif")
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/animation_example.gif)
</hfoption>
</hfoptions>
@@ -259,6 +259,50 @@ pip install git+https://github.com/huggingface/peft.git
**Inference**
The inference is the same as if you train a regular LoRA 🤗
## Conducting EDM-style training
It's now possible to perform EDM-style training as proposed in [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364).
simply set:
```diff
+ --do_edm_style_training \
```
Other SDXL-like models that use the EDM formulation, such as [playgroundai/playground-v2.5-1024px-aesthetic](https://huggingface.co/playgroundai/playground-v2.5-1024px-aesthetic), can also be DreamBooth'd with the script. Below is an example command:
```bash
accelerate launch train_dreambooth_lora_sdxl_advanced.py \
--pretrained_model_name_or_path="playgroundai/playground-v2.5-1024px-aesthetic" \
--dataset_name="linoyts/3d_icon" \
--instance_prompt="3d icon in the style of TOK" \
--validation_prompt="a TOK icon of an astronaut riding a horse, in the style of TOK" \
--output_dir="3d-icon-SDXL-LoRA" \
--do_edm_style_training \
--caption_column="prompt" \
--mixed_precision="bf16" \
--resolution=1024 \
--train_batch_size=3 \
--repeats=1 \
--report_to="wandb"\
--gradient_accumulation_steps=1 \
--gradient_checkpointing \
--learning_rate=1.0 \
--text_encoder_lr=1.0 \
--optimizer="prodigy"\
--train_text_encoder_ti\
--train_text_encoder_ti_frac=0.5\
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--rank=8 \
--max_train_steps=1000 \
--checkpointing_steps=2000 \
--seed="0" \
--push_to_hub
```
> [!CAUTION]
> Min-SNR gamma is not supported with the EDM-style training yet. When training with the PlaygroundAI model, it's recommended to not pass any "variant".
### Tips and Tricks
Check out [these recommended practices](https://huggingface.co/blog/sdxl_lora_advanced_script#additional-good-practices)
@@ -70,7 +70,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -14,9 +14,11 @@
# See the License for the specific language governing permissions and
import argparse
import contextlib
import gc
import hashlib
import itertools
import json
import logging
import math
import os
@@ -37,7 +39,7 @@ import transformers
from accelerate import Accelerator
from accelerate.logging import get_logger
from accelerate.utils import DistributedDataParallelKwargs, ProjectConfiguration, set_seed
from huggingface_hub import create_repo, upload_folder
from huggingface_hub import create_repo, hf_hub_download, upload_folder
from packaging import version
from peft import LoraConfig, set_peft_model_state_dict
from peft.utils import get_peft_model_state_dict
@@ -55,6 +57,8 @@ from diffusers import (
AutoencoderKL,
DDPMScheduler,
DPMSolverMultistepScheduler,
EDMEulerScheduler,
EulerDiscreteScheduler,
StableDiffusionXLPipeline,
UNet2DConditionModel,
)
@@ -74,11 +78,25 @@ from diffusers.utils.torch_utils import is_compiled_module
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
def determine_scheduler_type(pretrained_model_name_or_path, revision):
model_index_filename = "model_index.json"
if os.path.isdir(pretrained_model_name_or_path):
model_index = os.path.join(pretrained_model_name_or_path, model_index_filename)
else:
model_index = hf_hub_download(
repo_id=pretrained_model_name_or_path, filename=model_index_filename, revision=revision
)
with open(model_index, "r") as f:
scheduler_type = json.load(f)["scheduler"][1]
return scheduler_type
def save_model_card(
repo_id: str,
use_dora: bool,
@@ -370,6 +388,11 @@ def parse_args(input_args=None):
" `args.validation_prompt` multiple times: `args.num_validation_images`."
),
)
parser.add_argument(
"--do_edm_style_training",
action="store_true",
help="Flag to conduct training using the EDM formulation as introduced in https://arxiv.org/abs/2206.00364.",
)
parser.add_argument(
"--with_prior_preservation",
default=False,
@@ -1117,6 +1140,8 @@ def main(args):
"You cannot use both --report_to=wandb and --hub_token due to a security risk of exposing your token."
" Please use `huggingface-cli login` to authenticate with the Hub."
)
if args.do_edm_style_training and args.snr_gamma is not None:
raise ValueError("Min-SNR formulation is not supported when conducting EDM-style training.")
logging_dir = Path(args.output_dir, args.logging_dir)
@@ -1234,7 +1259,19 @@ def main(args):
)
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
scheduler_type = determine_scheduler_type(args.pretrained_model_name_or_path, args.revision)
if "EDM" in scheduler_type:
args.do_edm_style_training = True
noise_scheduler = EDMEulerScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
logger.info("Performing EDM-style training!")
elif args.do_edm_style_training:
noise_scheduler = EulerDiscreteScheduler.from_pretrained(
args.pretrained_model_name_or_path, subfolder="scheduler"
)
logger.info("Performing EDM-style training!")
else:
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder_one = text_encoder_cls_one.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
@@ -1252,7 +1289,12 @@ def main(args):
revision=args.revision,
variant=args.variant,
)
vae_scaling_factor = vae.config.scaling_factor
latents_mean = latents_std = None
if hasattr(vae.config, "latents_mean") and vae.config.latents_mean is not None:
latents_mean = torch.tensor(vae.config.latents_mean).view(1, 4, 1, 1)
if hasattr(vae.config, "latents_std") and vae.config.latents_std is not None:
latents_std = torch.tensor(vae.config.latents_std).view(1, 4, 1, 1)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
@@ -1790,6 +1832,19 @@ def main(args):
disable=not accelerator.is_local_main_process,
)
def get_sigmas(timesteps, n_dim=4, dtype=torch.float32):
# TODO: revisit other sampling algorithms
sigmas = noise_scheduler.sigmas.to(device=accelerator.device, dtype=dtype)
schedule_timesteps = noise_scheduler.timesteps.to(accelerator.device)
timesteps = timesteps.to(accelerator.device)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < n_dim:
sigma = sigma.unsqueeze(-1)
return sigma
if args.train_text_encoder:
num_train_epochs_text_encoder = int(args.train_text_encoder_frac * args.num_train_epochs)
elif args.train_text_encoder_ti: # args.train_text_encoder_ti
@@ -1841,9 +1896,15 @@ def main(args):
pixel_values = batch["pixel_values"].to(dtype=vae.dtype)
model_input = vae.encode(pixel_values).latent_dist.sample()
model_input = model_input * vae_scaling_factor
if args.pretrained_vae_model_name_or_path is None:
model_input = model_input.to(weight_dtype)
if latents_mean is None and latents_std is None:
model_input = model_input * vae.config.scaling_factor
if args.pretrained_vae_model_name_or_path is None:
model_input = model_input.to(weight_dtype)
else:
latents_mean = latents_mean.to(device=model_input.device, dtype=model_input.dtype)
latents_std = latents_std.to(device=model_input.device, dtype=model_input.dtype)
model_input = (model_input - latents_mean) * vae.config.scaling_factor / latents_std
model_input = model_input.to(dtype=weight_dtype)
# Sample noise that we'll add to the latents
noise = torch.randn_like(model_input)
@@ -1854,15 +1915,32 @@ def main(args):
)
bsz = model_input.shape[0]
# Sample a random timestep for each image
timesteps = torch.randint(
0, noise_scheduler.config.num_train_timesteps, (bsz,), device=model_input.device
)
timesteps = timesteps.long()
if not args.do_edm_style_training:
timesteps = torch.randint(
0, noise_scheduler.config.num_train_timesteps, (bsz,), device=model_input.device
)
timesteps = timesteps.long()
else:
# in EDM formulation, the model is conditioned on the pre-conditioned noise levels
# instead of discrete timesteps, so here we sample indices to get the noise levels
# from `scheduler.timesteps`
indices = torch.randint(0, noise_scheduler.config.num_train_timesteps, (bsz,))
timesteps = noise_scheduler.timesteps[indices].to(device=model_input.device)
# Add noise to the model input according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_model_input = noise_scheduler.add_noise(model_input, noise, timesteps)
# For EDM-style training, we first obtain the sigmas based on the continuous timesteps.
# We then precondition the final model inputs based on these sigmas instead of the timesteps.
# Follow: Section 5 of https://arxiv.org/abs/2206.00364.
if args.do_edm_style_training:
sigmas = get_sigmas(timesteps, len(noisy_model_input.shape), noisy_model_input.dtype)
if "EDM" in scheduler_type:
inp_noisy_latents = noise_scheduler.precondition_inputs(noisy_model_input, sigmas)
else:
inp_noisy_latents = noisy_model_input / ((sigmas**2 + 1) ** 0.5)
# time ids
add_time_ids = torch.cat(
@@ -1888,7 +1966,7 @@ def main(args):
}
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat_text_embeds, 1, 1)
model_pred = unet(
noisy_model_input,
inp_noisy_latents if args.do_edm_style_training else noisy_model_input,
timesteps,
prompt_embeds_input,
added_cond_kwargs=unet_added_conditions,
@@ -1906,14 +1984,42 @@ def main(args):
)
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat_text_embeds, 1, 1)
model_pred = unet(
noisy_model_input, timesteps, prompt_embeds_input, added_cond_kwargs=unet_added_conditions
inp_noisy_latents if args.do_edm_style_training else noisy_model_input,
timesteps,
prompt_embeds_input,
added_cond_kwargs=unet_added_conditions,
).sample
weighting = None
if args.do_edm_style_training:
# Similar to the input preconditioning, the model predictions are also preconditioned
# on noised model inputs (before preconditioning) and the sigmas.
# Follow: Section 5 of https://arxiv.org/abs/2206.00364.
if "EDM" in scheduler_type:
model_pred = noise_scheduler.precondition_outputs(noisy_model_input, model_pred, sigmas)
else:
if noise_scheduler.config.prediction_type == "epsilon":
model_pred = model_pred * (-sigmas) + noisy_model_input
elif noise_scheduler.config.prediction_type == "v_prediction":
model_pred = model_pred * (-sigmas / (sigmas**2 + 1) ** 0.5) + (
noisy_model_input / (sigmas**2 + 1)
)
# We are not doing weighting here because it tends result in numerical problems.
# See: https://github.com/huggingface/diffusers/pull/7126#issuecomment-1968523051
# There might be other alternatives for weighting as well:
# https://github.com/huggingface/diffusers/pull/7126#discussion_r1505404686
if "EDM" not in scheduler_type:
weighting = (sigmas**-2.0).float()
# Get the target for loss depending on the prediction type
if noise_scheduler.config.prediction_type == "epsilon":
target = noise
target = model_input if args.do_edm_style_training else noise
elif noise_scheduler.config.prediction_type == "v_prediction":
target = noise_scheduler.get_velocity(model_input, noise, timesteps)
target = (
model_input
if args.do_edm_style_training
else noise_scheduler.get_velocity(model_input, noise, timesteps)
)
else:
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
@@ -1923,10 +2029,28 @@ def main(args):
target, target_prior = torch.chunk(target, 2, dim=0)
# Compute prior loss
prior_loss = F.mse_loss(model_pred_prior.float(), target_prior.float(), reduction="mean")
if weighting is not None:
prior_loss = torch.mean(
(weighting.float() * (model_pred_prior.float() - target_prior.float()) ** 2).reshape(
target_prior.shape[0], -1
),
1,
)
prior_loss = prior_loss.mean()
else:
prior_loss = F.mse_loss(model_pred_prior.float(), target_prior.float(), reduction="mean")
if args.snr_gamma is None:
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
if weighting is not None:
loss = torch.mean(
(weighting.float() * (model_pred.float() - target.float()) ** 2).reshape(
target.shape[0], -1
),
1,
)
loss = loss.mean()
else:
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
else:
# Compute loss-weights as per Section 3.4 of https://arxiv.org/abs/2303.09556.
# Since we predict the noise instead of x_0, the original formulation is slightly changed.
@@ -2049,17 +2173,18 @@ def main(args):
# We train on the simplified learning objective. If we were previously predicting a variance, we need the scheduler to ignore it
scheduler_args = {}
if "variance_type" in pipeline.scheduler.config:
variance_type = pipeline.scheduler.config.variance_type
if not args.do_edm_style_training:
if "variance_type" in pipeline.scheduler.config:
variance_type = pipeline.scheduler.config.variance_type
if variance_type in ["learned", "learned_range"]:
variance_type = "fixed_small"
if variance_type in ["learned", "learned_range"]:
variance_type = "fixed_small"
scheduler_args["variance_type"] = variance_type
scheduler_args["variance_type"] = variance_type
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(
pipeline.scheduler.config, **scheduler_args
)
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(
pipeline.scheduler.config, **scheduler_args
)
pipeline = pipeline.to(accelerator.device)
pipeline.set_progress_bar_config(disable=True)
@@ -2067,8 +2192,13 @@ def main(args):
# run inference
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
pipeline_args = {"prompt": args.validation_prompt}
inference_ctx = (
contextlib.nullcontext()
if "playground" in args.pretrained_model_name_or_path
else torch.cuda.amp.autocast()
)
with torch.cuda.amp.autocast():
with inference_ctx:
images = [
pipeline(**pipeline_args, generator=generator).images[0]
for _ in range(args.num_validation_images)
@@ -2144,15 +2274,18 @@ def main(args):
# We train on the simplified learning objective. If we were previously predicting a variance, we need the scheduler to ignore it
scheduler_args = {}
if "variance_type" in pipeline.scheduler.config:
variance_type = pipeline.scheduler.config.variance_type
if not args.do_edm_style_training:
if "variance_type" in pipeline.scheduler.config:
variance_type = pipeline.scheduler.config.variance_type
if variance_type in ["learned", "learned_range"]:
variance_type = "fixed_small"
if variance_type in ["learned", "learned_range"]:
variance_type = "fixed_small"
scheduler_args["variance_type"] = variance_type
scheduler_args["variance_type"] = variance_type
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config, **scheduler_args)
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(
pipeline.scheduler.config, **scheduler_args
)
# load attention processors
pipeline.load_lora_weights(args.output_dir)
@@ -513,9 +513,7 @@ class LCMSchedulerWithTimestamp(SchedulerMixin, ConfigMixin):
there is no previous alpha. When this option is `True` the previous alpha product is fixed to `1`,
otherwise it uses the alpha value at step 0.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, defaults to `epsilon`, *optional*):
Prediction type of the scheduler function; can be `epsilon` (predicts the noise of the diffusion process),
`sample` (directly predicts the noisy sample`) or `v_prediction` (see section 2.4 of [Imagen
@@ -418,9 +418,7 @@ class LCMScheduler(SchedulerMixin, ConfigMixin):
there is no previous alpha. When this option is `True` the previous alpha product is fixed to `1`,
otherwise it uses the alpha value at step 0.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, defaults to `epsilon`, *optional*):
Prediction type of the scheduler function; can be `epsilon` (predicts the noise of the diffusion process),
`sample` (directly predicts the noisy sample`) or `v_prediction` (see section 2.4 of [Imagen
@@ -40,7 +40,7 @@ from diffusers.utils import BaseOutput, check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
class MarigoldDepthOutput(BaseOutput):
+1 -3
View File
@@ -171,9 +171,7 @@ class UFOGenScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -72,7 +72,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -65,7 +65,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -78,7 +78,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -71,7 +71,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -77,7 +77,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
+1 -1
View File
@@ -60,7 +60,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
+1 -1
View File
@@ -60,7 +60,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = logging.getLogger(__name__)
+1 -1
View File
@@ -61,7 +61,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -63,7 +63,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -1178,7 +1178,7 @@ def main(args):
grads_text_encoder = text_encoder.get_input_embeddings().weight.grad
# Get the index for tokens that we want to zero the grads for
index_grads_to_zero = torch.arange(len(tokenizer)) != modifier_token_id[0]
for i in range(len(modifier_token_id[1:])):
for i in range(1, len(modifier_token_id)):
index_grads_to_zero = index_grads_to_zero & (
torch.arange(len(tokenizer)) != modifier_token_id[i]
)
+1 -1
View File
@@ -63,7 +63,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
+1 -1
View File
@@ -35,7 +35,7 @@ from diffusers.utils import check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
# Cache compiled models across invocations of this script.
cc.initialize_cache(os.path.expanduser("~/.cache/jax/compilation_cache"))
+1 -1
View File
@@ -70,7 +70,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -75,7 +75,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -53,7 +53,7 @@ from diffusers.utils.torch_utils import is_compiled_module
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -59,7 +59,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -52,7 +52,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -51,7 +51,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -637,7 +637,7 @@ def main(args):
generator=generator,
batch_size=args.eval_batch_size,
num_inference_steps=args.ddpm_num_inference_steps,
output_type="numpy",
output_type="np",
).images
if args.use_ema:
@@ -60,7 +60,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -56,7 +56,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -49,7 +49,7 @@ from diffusers.utils import check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = logging.getLogger(__name__)
@@ -52,7 +52,7 @@ from diffusers.utils.torch_utils import is_compiled_module
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -64,7 +64,7 @@ from diffusers.utils.torch_utils import is_compiled_module
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -425,6 +425,11 @@ def parse_args(input_args=None):
default=4,
help=("The dimension of the LoRA update matrices."),
)
parser.add_argument(
"--debug_loss",
action="store_true",
help="debug loss for each image, if filenames are awailable in the dataset",
)
if input_args is not None:
args = parser.parse_args(input_args)
@@ -603,6 +608,7 @@ def main(args):
# Move unet, vae and text_encoder to device and cast to weight_dtype
# The VAE is in float32 to avoid NaN losses.
unet.to(accelerator.device, dtype=weight_dtype)
if args.pretrained_vae_model_name_or_path is None:
vae.to(accelerator.device, dtype=torch.float32)
else:
@@ -890,13 +896,17 @@ def main(args):
tokens_one, tokens_two = tokenize_captions(examples)
examples["input_ids_one"] = tokens_one
examples["input_ids_two"] = tokens_two
if args.debug_loss:
fnames = [os.path.basename(image.filename) for image in examples[image_column] if image.filename]
if fnames:
examples["filenames"] = fnames
return examples
with accelerator.main_process_first():
if args.max_train_samples is not None:
dataset["train"] = dataset["train"].shuffle(seed=args.seed).select(range(args.max_train_samples))
# Set the training transforms
train_dataset = dataset["train"].with_transform(preprocess_train)
train_dataset = dataset["train"].with_transform(preprocess_train, output_all_columns=True)
def collate_fn(examples):
pixel_values = torch.stack([example["pixel_values"] for example in examples])
@@ -905,7 +915,7 @@ def main(args):
crop_top_lefts = [example["crop_top_lefts"] for example in examples]
input_ids_one = torch.stack([example["input_ids_one"] for example in examples])
input_ids_two = torch.stack([example["input_ids_two"] for example in examples])
return {
result = {
"pixel_values": pixel_values,
"input_ids_one": input_ids_one,
"input_ids_two": input_ids_two,
@@ -913,6 +923,11 @@ def main(args):
"crop_top_lefts": crop_top_lefts,
}
filenames = [example["filenames"] for example in examples if "filenames" in example]
if filenames:
result["filenames"] = filenames
return result
# DataLoaders creation:
train_dataloader = torch.utils.data.DataLoader(
train_dataset,
@@ -1105,7 +1120,9 @@ def main(args):
loss = F.mse_loss(model_pred.float(), target.float(), reduction="none")
loss = loss.mean(dim=list(range(1, len(loss.shape)))) * mse_loss_weights
loss = loss.mean()
if args.debug_loss and "filenames" in batch:
for fname in batch["filenames"]:
accelerator.log({"loss_for_" + fname: loss}, step=global_step)
# Gather the losses across all processes for logging (if we use distributed training).
avg_loss = accelerator.gather(loss.repeat(args.train_batch_size)).mean()
train_loss += avg_loss.item() / args.gradient_accumulation_steps
@@ -54,7 +54,7 @@ from diffusers.utils.torch_utils import is_compiled_module
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -80,7 +80,7 @@ else:
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -56,7 +56,7 @@ else:
# ------------------------------------------------------------------------------
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = logging.getLogger(__name__)
@@ -76,7 +76,7 @@ else:
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__)
@@ -29,7 +29,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -648,7 +648,7 @@ def main(args):
generator=generator,
batch_size=args.eval_batch_size,
num_inference_steps=args.ddpm_num_inference_steps,
output_type="numpy",
output_type="np",
).images
if args.use_ema:
@@ -50,7 +50,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -51,7 +51,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.27.0.dev0")
check_min_version("0.28.0.dev0")
logger = get_logger(__name__, log_level="INFO")
+139
View File
@@ -0,0 +1,139 @@
import argparse
import json
import os
from datetime import date
from pathlib import Path
from slack_sdk import WebClient
from tabulate import tabulate
MAX_LEN_MESSAGE = 2900 # slack endpoint has a limit of 3001 characters
parser = argparse.ArgumentParser()
parser.add_argument("--slack_channel_name", default="diffusers-ci-nightly")
def main(slack_channel_name=None):
failed = []
passed = []
group_info = []
total_num_failed = 0
empty_file = False or len(list(Path().glob("*.log"))) == 0
total_empty_files = []
for log in Path().glob("*.log"):
section_num_failed = 0
i = 0
with open(log) as f:
for line in f:
line = json.loads(line)
i += 1
if line.get("nodeid", "") != "":
test = line["nodeid"]
if line.get("duration", None) is not None:
duration = f'{line["duration"]:.4f}'
if line.get("outcome", "") == "failed":
section_num_failed += 1
failed.append([test, duration, log.name.split("_")[0]])
total_num_failed += 1
else:
passed.append([test, duration, log.name.split("_")[0]])
empty_file = i == 0
group_info.append([str(log), section_num_failed, failed])
total_empty_files.append(empty_file)
os.remove(log)
failed = []
text = (
"🌞 There were no failures!"
if not any(total_empty_files)
else "Something went wrong there is at least one empty file - please check GH action results."
)
no_error_payload = {
"type": "section",
"text": {
"type": "plain_text",
"text": text,
"emoji": True,
},
}
message = ""
payload = [
{
"type": "header",
"text": {
"type": "plain_text",
"text": "🤗 Results of the Diffusers scheduled nightly tests.",
},
},
]
if total_num_failed > 0:
for i, (name, num_failed, failed_tests) in enumerate(group_info):
if num_failed > 0:
if num_failed == 1:
message += f"*{name}: {num_failed} failed test*\n"
else:
message += f"*{name}: {num_failed} failed tests*\n"
failed_table = []
for test in failed_tests:
failed_table.append(test[0].split("::"))
failed_table = tabulate(
failed_table,
headers=["Test Location", "Test Case", "Test Name"],
showindex="always",
tablefmt="grid",
maxcolwidths=[12, 12, 12],
)
message += "\n```\n" + failed_table + "\n```"
if total_empty_files[i]:
message += f"\n*{name}: Warning! Empty file - please check the GitHub action job *\n"
print(f"### {message}")
else:
payload.append(no_error_payload)
if len(message) > MAX_LEN_MESSAGE:
print(f"Truncating long message from {len(message)} to {MAX_LEN_MESSAGE}")
message = message[:MAX_LEN_MESSAGE] + "..."
if len(message) != 0:
md_report = {
"type": "section",
"text": {"type": "mrkdwn", "text": message},
}
payload.append(md_report)
action_button = {
"type": "section",
"text": {"type": "mrkdwn", "text": "*For more details:*"},
"accessory": {
"type": "button",
"text": {"type": "plain_text", "text": "Check Action results", "emoji": True},
"url": f"https://github.com/huggingface/diffusers/actions/runs/{os.environ['GITHUB_RUN_ID']}",
},
}
payload.append(action_button)
date_report = {
"type": "context",
"elements": [
{
"type": "plain_text",
"text": f"Nightly test results for {date.today()}",
},
],
}
payload.append(date_report)
print(payload)
client = WebClient(token=os.environ.get("SLACK_API_TOKEN"))
client.chat_postMessage(channel=f"#{slack_channel_name}", text=message, blocks=payload)
if __name__ == "__main__":
args = parser.parse_args()
main(args.slack_channel_name)
+1 -1
View File
@@ -249,7 +249,7 @@ version_range_max = max(sys.version_info[1], 10) + 1
setup(
name="diffusers",
version="0.27.0.dev0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
version="0.28.0.dev0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
description="State-of-the-art diffusion in PyTorch and JAX.",
long_description=open("README.md", "r", encoding="utf-8").read(),
long_description_content_type="text/markdown",
+5 -1
View File
@@ -1,4 +1,4 @@
__version__ = "0.27.0.dev0"
__version__ = "0.28.0.dev0"
from typing import TYPE_CHECKING
@@ -253,6 +253,8 @@ else:
"LatentConsistencyModelImg2ImgPipeline",
"LatentConsistencyModelPipeline",
"LDMTextToImagePipeline",
"LEditsPPPipelineStableDiffusion",
"LEditsPPPipelineStableDiffusionXL",
"MusicLDMPipeline",
"PaintByExamplePipeline",
"PIAPipeline",
@@ -623,6 +625,8 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
LatentConsistencyModelImg2ImgPipeline,
LatentConsistencyModelPipeline,
LDMTextToImagePipeline,
LEditsPPPipelineStableDiffusion,
LEditsPPPipelineStableDiffusionXL,
MusicLDMPipeline,
PaintByExamplePipeline,
PIAPipeline,
+1 -1
View File
@@ -293,7 +293,7 @@ class BasicTransformerBlock(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
# Notice that normalization is always applied before the real computation in the following blocks.
# 0. Self-Attention
+24 -2
View File
@@ -767,7 +767,18 @@ class AttnProcessor:
query = attn.to_q(hidden_states)
if encoder_hidden_states is None:
encoder_hidden_states = hidden_states
# encoder_hidden_states = hidden_states
batch, seq, dim = hidden_states.shape
height = width = seq**0.5
# reshape to (batch, height, width, dim)
encoder_hidden_states = hidden_states.view(batch, height, width, dim)
# reshape to (batch, dim, height, width)
encoder_hidden_states = encoder_hidden_states.permute(0, 3, 1, 2)
encoder_hidden_states = torch.nn.functional.avg_pool2d(hidden_states, kernel_size=4)
# reshape to (batch, dim, seq)
encoder_hidden_states = encoder_hidden_states.view(batch, dim, -1)
# reshape to (batch, seq, dim)
encoder_hidden_states = encoder_hidden_states.permute(0, 2, 1)
elif attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
@@ -1259,7 +1270,18 @@ class AttnProcessor2_0:
query = attn.to_q(hidden_states)
if encoder_hidden_states is None:
encoder_hidden_states = hidden_states
# encoder_hidden_states = hidden_states
batch, seq, dim = hidden_states.shape
height = width = seq**0.5
# reshape to (batch, height, width, dim)
encoder_hidden_states = hidden_states.view(batch, height, width, dim)
# reshape to (batch, dim, height, width)
encoder_hidden_states = encoder_hidden_states.permute(0, 3, 1, 2)
encoder_hidden_states = torch.nn.functional.avg_pool2d(hidden_states, kernel_size=4)
# reshape to (batch, dim, seq)
encoder_hidden_states = encoder_hidden_states.view(batch, dim, -1)
# reshape to (batch, seq, dim)
encoder_hidden_states = encoder_hidden_states.permute(0, 2, 1)
elif attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
@@ -308,7 +308,7 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
"""
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
# ensure attention_mask is a bias, and give it a singleton query_tokens dimension.
# we may have done this conversion already, e.g. if we came here via UNet2DConditionModel#forward.
# we can tell by counting dims; if ndim == 2: it's a mask rather than a bias.
+9 -9
View File
@@ -846,7 +846,7 @@ class UNetMidBlock2DCrossAttn(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
hidden_states = self.resnets[0](hidden_states, temb)
for attn, resnet in zip(self.attentions, self.resnets[1:]):
@@ -986,7 +986,7 @@ class UNetMidBlock2DSimpleCrossAttn(nn.Module):
) -> torch.FloatTensor:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
if attention_mask is None:
# if encoder_hidden_states is defined: we are doing cross-attn, so we should use cross-attn mask.
@@ -1116,7 +1116,7 @@ class AttnDownBlock2D(nn.Module):
) -> Tuple[torch.FloatTensor, Tuple[torch.FloatTensor, ...]]:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
output_states = ()
@@ -1241,7 +1241,7 @@ class CrossAttnDownBlock2D(nn.Module):
) -> Tuple[torch.FloatTensor, Tuple[torch.FloatTensor, ...]]:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
output_states = ()
@@ -1986,7 +1986,7 @@ class SimpleCrossAttnDownBlock2D(nn.Module):
) -> Tuple[torch.FloatTensor, Tuple[torch.FloatTensor, ...]]:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
output_states = ()
@@ -2201,7 +2201,7 @@ class KCrossAttnDownBlock2D(nn.Module):
) -> Tuple[torch.FloatTensor, Tuple[torch.FloatTensor, ...]]:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
output_states = ()
@@ -2483,7 +2483,7 @@ class CrossAttnUpBlock2D(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
is_freeu_enabled = (
getattr(self, "s1", None)
@@ -3312,7 +3312,7 @@ class SimpleCrossAttnUpBlock2D(nn.Module):
) -> torch.FloatTensor:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
if attention_mask is None:
# if encoder_hidden_states is defined: we are doing cross-attn, so we should use cross-attn mask.
@@ -3694,7 +3694,7 @@ class KAttentionBlock(nn.Module):
) -> torch.FloatTensor:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
# 1. Self-Attention
if self.add_self_attention:
@@ -80,7 +80,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin, UNet2DConditionLoadersMixin,
in_channels (`int`, *optional*, defaults to 4): Number of channels in the input sample.
out_channels (`int`, *optional*, defaults to 4): Number of channels in the output.
center_input_sample (`bool`, *optional*, defaults to `False`): Whether to center the input sample.
flip_sin_to_cos (`bool`, *optional*, defaults to `False`):
flip_sin_to_cos (`bool`, *optional*, defaults to `True`):
Whether to flip the sin to cos in the time embedding.
freq_shift (`int`, *optional*, defaults to 0): The frequency shift to apply to the time embedding.
down_block_types (`Tuple[str]`, *optional*, defaults to `("CrossAttnDownBlock2D", "CrossAttnDownBlock2D", "CrossAttnDownBlock2D", "DownBlock2D")`):
@@ -109,7 +109,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin, UNet2DConditionLoadersMixin,
The number of transformer blocks of type [`~models.attention.BasicTransformerBlock`]. Only relevant for
[`~models.unet_2d_blocks.CrossAttnDownBlock2D`], [`~models.unet_2d_blocks.CrossAttnUpBlock2D`],
[`~models.unet_2d_blocks.UNetMidBlock2DCrossAttn`].
reverse_transformer_layers_per_block : (`Tuple[Tuple]`, *optional*, defaults to None):
reverse_transformer_layers_per_block : (`Tuple[Tuple]`, *optional*, defaults to None):
The number of transformer blocks of type [`~models.attention.BasicTransformerBlock`], in the upsampling
blocks of the U-Net. Only relevant if `transformer_layers_per_block` is of type `Tuple[Tuple]` and for
[`~models.unet_2d_blocks.CrossAttnDownBlock2D`], [`~models.unet_2d_blocks.CrossAttnUpBlock2D`],
@@ -147,9 +147,9 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin, UNet2DConditionLoadersMixin,
The second activation function to use in timestep embedding. Choose from `silu`, `mish` and `gelu`.
time_cond_proj_dim (`int`, *optional*, defaults to `None`):
The dimension of `cond_proj` layer in the timestep embedding.
conv_in_kernel (`int`, *optional*, default to `3`): The kernel size of `conv_in` layer. conv_out_kernel (`int`,
*optional*, default to `3`): The kernel size of `conv_out` layer. projection_class_embeddings_input_dim (`int`,
*optional*): The dimension of the `class_labels` input when
conv_in_kernel (`int`, *optional*, default to `3`): The kernel size of `conv_in` layer.
conv_out_kernel (`int`, *optional*, default to `3`): The kernel size of `conv_out` layer.
projection_class_embeddings_input_dim (`int`, *optional*): The dimension of the `class_labels` input when
`class_embed_type="projection"`. Required when `class_embed_type="projection"`.
class_embeddings_concat (`bool`, *optional*, defaults to `False`): Whether to concatenate the time
embeddings with the class embeddings.
+3 -3
View File
@@ -1183,7 +1183,7 @@ class CrossAttnDownBlockMotion(nn.Module):
):
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
output_states = ()
@@ -1367,7 +1367,7 @@ class CrossAttnUpBlockMotion(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
is_freeu_enabled = (
getattr(self, "s1", None)
@@ -1707,7 +1707,7 @@ class UNetMidBlockCrossAttnMotion(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
hidden_states = self.resnets[0](hidden_states, temb)
+13
View File
@@ -23,6 +23,7 @@ _import_structure = {
"controlnet_xs": [],
"deprecated": [],
"latent_diffusion": [],
"ledits_pp": [],
"stable_diffusion": [],
"stable_diffusion_xl": [],
}
@@ -171,6 +172,12 @@ else:
"LatentConsistencyModelPipeline",
]
_import_structure["latent_diffusion"].extend(["LDMTextToImagePipeline"])
_import_structure["ledits_pp"].extend(
[
"LEditsPPPipelineStableDiffusion",
"LEditsPPPipelineStableDiffusionXL",
]
)
_import_structure["musicldm"] = ["MusicLDMPipeline"]
_import_structure["paint_by_example"] = ["PaintByExamplePipeline"]
_import_structure["pia"] = ["PIAPipeline"]
@@ -424,6 +431,12 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
LatentConsistencyModelPipeline,
)
from .latent_diffusion import LDMTextToImagePipeline
from .ledits_pp import (
LEditsPPDiffusionPipelineOutput,
LEditsPPInversionPipelineOutput,
LEditsPPPipelineStableDiffusion,
LEditsPPPipelineStableDiffusionXL,
)
from .musicldm import MusicLDMPipeline
from .paint_by_example import PaintByExamplePipeline
from .pia import PIAPipeline
@@ -127,7 +127,7 @@ class AmusedImg2ImgPipeline(DiffusionPipeline):
on the amount of noise initially added. When `strength` is 1, added noise is maximum and the denoising
process runs for the full number of iterations specified in `num_inference_steps`. A value of 1
essentially ignores `image`.
num_inference_steps (`int`, *optional*, defaults to 16):
num_inference_steps (`int`, *optional*, defaults to 12):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
guidance_scale (`float`, *optional*, defaults to 10.0):
@@ -191,7 +191,7 @@ class AmusedImg2ImgPipeline(DiffusionPipeline):
negative_prompt_embeds is None and negative_encoder_hidden_states is not None
):
raise ValueError(
"pass either both `negatve_prompt_embeds` and `negative_encoder_hidden_states` or neither"
"pass either both `negative_prompt_embeds` and `negative_encoder_hidden_states` or neither"
)
if (prompt is None and prompt_embeds is None) or (prompt is not None and prompt_embeds is not None):
@@ -824,20 +824,22 @@ class StableDiffusionControlNetPipeline(
return latents
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -869,20 +869,22 @@ class StableDiffusionXLControlNetPipeline(
self.vae.decoder.mid_block.to(dtype)
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -133,7 +133,7 @@ class SpectrogramDiffusionPipeline(DiffusionPipeline):
generator: Optional[torch.Generator] = None,
num_inference_steps: int = 100,
return_dict: bool = True,
output_type: str = "numpy",
output_type: str = "np",
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: int = 1,
) -> Union[AudioPipelineOutput, Tuple]:
@@ -157,7 +157,7 @@ class SpectrogramDiffusionPipeline(DiffusionPipeline):
expense of slower inference.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~pipelines.AudioPipelineOutput`] instead of a plain tuple.
output_type (`str`, *optional*, defaults to `"numpy"`):
output_type (`str`, *optional*, defaults to `"np"`):
The output format of the generated audio.
callback (`Callable`, *optional*):
A function that calls every `callback_steps` steps during inference. The function is called with the
@@ -249,16 +249,16 @@ class SpectrogramDiffusionPipeline(DiffusionPipeline):
logger.info("Generated segment", i)
if output_type == "numpy" and not is_onnx_available():
if output_type == "np" and not is_onnx_available():
raise ValueError(
"Cannot return output in 'np' format if ONNX is not available. Make sure to have ONNX installed or set 'output_type' to 'mel'."
)
elif output_type == "numpy" and self.melgan is None:
elif output_type == "np" and self.melgan is None:
raise ValueError(
"Cannot return output in 'np' format if melgan component is not defined. Make sure to define `self.melgan` or set 'output_type' to 'mel'."
)
if output_type == "numpy":
if output_type == "np":
output = self.melgan(input_features=full_pred_mel.astype(np.float32))
else:
output = full_pred_mel
@@ -2004,7 +2004,7 @@ class CrossAttnUpBlockFlat(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
is_freeu_enabled = (
getattr(self, "s1", None)
@@ -2338,7 +2338,7 @@ class UNetMidBlockFlatCrossAttn(nn.Module):
) -> torch.FloatTensor:
if cross_attention_kwargs is not None:
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
hidden_states = self.resnets[0](hidden_states, temb)
for attn, resnet in zip(self.attentions, self.resnets[1:]):
@@ -2479,7 +2479,7 @@ class UNetMidBlockFlatSimpleCrossAttn(nn.Module):
) -> torch.FloatTensor:
cross_attention_kwargs = cross_attention_kwargs if cross_attention_kwargs is not None else {}
if cross_attention_kwargs.get("scale", None) is not None:
logger.warning("Passing `scale` to `cross_attention_kwargs` is depcrecated. `scale` will be ignored.")
logger.warning("Passing `scale` to `cross_attention_kwargs` is deprecated. `scale` will be ignored.")
if attention_mask is None:
# if encoder_hidden_states is defined: we are doing cross-attn, so we should use cross-attn mask.
@@ -548,20 +548,22 @@ class LatentConsistencyModelImg2ImgPipeline(
return latents
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -490,20 +490,22 @@ class LatentConsistencyModelPipeline(
latents = latents * self.scheduler.init_noise_sigma
return latents
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -0,0 +1,55 @@
from typing import TYPE_CHECKING
from ...utils import (
DIFFUSERS_SLOW_IMPORT,
OptionalDependencyNotAvailable,
_LazyModule,
get_objects_from_module,
is_torch_available,
is_transformers_available,
)
_dummy_objects = {}
_import_structure = {}
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_torch_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_torch_and_transformers_objects))
else:
_import_structure["pipeline_leditspp_stable_diffusion"] = ["LEditsPPPipelineStableDiffusion"]
_import_structure["pipeline_leditspp_stable_diffusion_xl"] = ["LEditsPPPipelineStableDiffusionXL"]
_import_structure["pipeline_output"] = ["LEditsPPDiffusionPipelineOutput", "LEditsPPDiffusionPipelineOutput"]
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_torch_and_transformers_objects import *
else:
from .pipeline_leditspp_stable_diffusion import (
LEditsPPDiffusionPipelineOutput,
LEditsPPInversionPipelineOutput,
LEditsPPPipelineStableDiffusion,
)
from .pipeline_leditspp_stable_diffusion_xl import LEditsPPPipelineStableDiffusionXL
else:
import sys
sys.modules[__name__] = _LazyModule(
__name__,
globals()["__file__"],
_import_structure,
module_spec=__spec__,
)
for name, value in _dummy_objects.items():
setattr(sys.modules[__name__], name, value)
File diff suppressed because it is too large Load Diff
File diff suppressed because it is too large Load Diff
@@ -0,0 +1,43 @@
from dataclasses import dataclass
from typing import List, Optional, Union
import numpy as np
import PIL.Image
from ...utils import BaseOutput
@dataclass
class LEditsPPDiffusionPipelineOutput(BaseOutput):
"""
Output class for LEdits++ Diffusion pipelines.
Args:
images (`List[PIL.Image.Image]` or `np.ndarray`)
List of denoised PIL images of length `batch_size` or NumPy array of shape `(batch_size, height, width,
num_channels)`.
nsfw_content_detected (`List[bool]`)
List indicating whether the corresponding generated image contains not-safe-for-work (nsfw) content or
`None` if safety checking could not be performed.
"""
images: Union[List[PIL.Image.Image], np.ndarray]
nsfw_content_detected: Optional[List[bool]]
@dataclass
class LEditsPPInversionPipelineOutput(BaseOutput):
"""
Output class for LEdits++ Diffusion pipelines.
Args:
input_images (`List[PIL.Image.Image]` or `np.ndarray`)
List of the cropped and resized input images as PIL images of length `batch_size` or NumPy array of shape `
(batch_size, height, width, num_channels)`.
vae_reconstruction_images (`List[PIL.Image.Image]` or `np.ndarray`)
List of VAE reconstruction of all input images as PIL images of length `batch_size` or NumPy array of shape `
(batch_size, height, width, num_channels)`.
"""
images: Union[List[PIL.Image.Image], np.ndarray]
vae_reconstruction_images: Union[List[PIL.Image.Image], np.ndarray]
@@ -669,20 +669,22 @@ class StableDiffusionPipeline(
return latents
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -767,20 +767,22 @@ class StableDiffusionImg2ImgPipeline(
return latents
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -909,20 +909,22 @@ class StableDiffusionInpaintPipeline(
return timesteps, num_inference_steps - t_start
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -1304,7 +1304,7 @@ class StableDiffusionDiffEditPipeline(
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: int = 1,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
clip_ckip: int = None,
clip_skip: int = None,
):
r"""
The call function to the pipeline for generation.
@@ -1426,7 +1426,7 @@ class StableDiffusionDiffEditPipeline(
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
lora_scale=text_encoder_lora_scale,
clip_skip=clip_ckip,
clip_skip=clip_skip,
)
# For classifier free guidance, we need to do two forward passes.
# Here we concatenate the unconditional and text embeddings into a single batch
@@ -644,20 +644,22 @@ class StableDiffusionLDM3DPipeline(
return latents
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -632,7 +632,7 @@ class StableDiffusionSAGPipeline(DiffusionPipeline, StableDiffusionMixin, Textua
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# and `sag_scale` is` `s` of equation (16)
# of the self-attentnion guidance paper: https://arxiv.org/pdf/2210.00939.pdf
# of the self-attention guidance paper: https://arxiv.org/pdf/2210.00939.pdf
# `sag_scale = 0` means no self-attention guidance
do_self_attention_guidance = sag_scale > 0.0
@@ -667,7 +667,7 @@ class StableDiffusionSAGPipeline(DiffusionPipeline, StableDiffusionMixin, Textua
if timesteps.dtype not in [torch.int16, torch.int32, torch.int64]:
raise ValueError(
f"{self.__class__.__name__} does not support using a scheduler of type {self.scheduler.__class__.__name__}. Please make sure to use one of 'DDIMScheduler, PNDMScheduler, DDPMScheduler, DEISMultistepScheduler, UniPCMultistepScheduler, DPMSolverMultistepScheduler, DPMSolverSinlgestepScheduler'."
f"{self.__class__.__name__} does not support using a scheduler of type {self.scheduler.__class__.__name__}. Please make sure to use one of 'DDIMScheduler, PNDMScheduler, DDPMScheduler, DEISMultistepScheduler, UniPCMultistepScheduler, DPMSolverMultistepScheduler, DPMSolverSinglestepScheduler'."
)
# 5. Prepare latent variables
@@ -723,7 +723,7 @@ class StableDiffusionSAGPipeline(DiffusionPipeline, StableDiffusionMixin, Textua
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
# perform self-attention guidance with the stored self-attentnion map
# perform self-attention guidance with the stored self-attention map
if do_self_attention_guidance:
# classifier-free guidance produces two chunks of attention map
# and we only use unconditional one according to equation (25)
@@ -740,20 +740,22 @@ class StableDiffusionXLPipeline(
self.vae.decoder.mid_block.to(dtype)
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -874,20 +874,22 @@ class StableDiffusionXLImg2ImgPipeline(
self.vae.decoder.mid_block.to(dtype)
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -1110,20 +1110,22 @@ class StableDiffusionXLInpaintPipeline(
self.vae.decoder.mid_block.to(dtype)
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -613,20 +613,22 @@ class StableDiffusionAdapterPipeline(DiffusionPipeline, StableDiffusionMixin):
return height, width
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -784,20 +784,22 @@ class StableDiffusionXLAdapterPipeline(
return height, width
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
def get_guidance_scale_embedding(
self, w: torch.Tensor, embedding_dim: int = 512, dtype: torch.dtype = torch.float32
) -> torch.FloatTensor:
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
w (`torch.Tensor`):
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Dimension of the embeddings to generate.
dtype (`torch.dtype`, *optional*, defaults to `torch.float32`):
Data type of the generated embeddings.
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
`torch.FloatTensor`: Embedding vectors with shape `(len(w), embedding_dim)`.
"""
assert len(w.shape) == 1
w = w * 1000.0
@@ -575,8 +575,8 @@ class TextToVideoZeroPipeline(DiffusionPipeline, StableDiffusionMixin, TextualIn
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for video
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random `generator`.
output_type (`str`, *optional*, defaults to `"numpy"`):
The output format of the generated video. Choose between `"latent"` and `"numpy"`.
output_type (`str`, *optional*, defaults to `"np"`):
The output format of the generated video. Choose between `"latent"` and `"np"`.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a
[`~pipelines.text_to_video_synthesis.pipeline_text_to_video_zero.TextToVideoPipelineOutput`] instead of
+1 -3
View File
@@ -157,9 +157,7 @@ class DDIMScheduler(SchedulerMixin, ConfigMixin):
there is no previous alpha. When this option is `True` the previous alpha product is fixed to `1`,
otherwise it uses the alpha value at step 0.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, defaults to `epsilon`, *optional*):
Prediction type of the scheduler function; can be `epsilon` (predicts the noise of the diffusion process),
`sample` (directly predicts the noisy sample`) or `v_prediction` (see section 2.4 of [Imagen
@@ -93,9 +93,7 @@ class FlaxDDIMScheduler(FlaxSchedulerMixin, ConfigMixin):
step there is no previous alpha. When this option is `True` the previous alpha product is fixed to `1`,
otherwise it uses the value of alpha at step 0.
steps_offset (`int`, default `0`):
an offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False`, to make the last step use step 0 for the previous alpha product, as done in
stable diffusion.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, default `epsilon`):
indicates whether the model predicts the noise (epsilon), or the samples. One of `epsilon`, `sample`.
`v-prediction` is not supported for this scheduler.
@@ -155,9 +155,7 @@ class DDIMInverseScheduler(SchedulerMixin, ConfigMixin):
there is no previous alpha. When this option is `True` the previous alpha product is fixed to 0, otherwise
it uses the alpha value at step `num_train_timesteps - 1`.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use `num_train_timesteps - 1` for the previous alpha
product.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, defaults to `epsilon`, *optional*):
Prediction type of the scheduler function; can be `epsilon` (predicts the noise of the diffusion process),
`sample` (directly predicts the noisy sample`) or `v_prediction` (see section 2.4 of [Imagen
@@ -159,9 +159,7 @@ class DDIMParallelScheduler(SchedulerMixin, ConfigMixin):
step there is no previous alpha. When this option is `True` the previous alpha product is fixed to `1`,
otherwise it uses the value of alpha at step 0.
steps_offset (`int`, default `0`):
an offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False`, to make the last step use step 0 for the previous alpha product, as done in
stable diffusion.
An offset added to the inference steps, as required by some model families.
prediction_type (`str`, default `epsilon`, optional):
prediction type of the scheduler function, one of `epsilon` (predicting the noise of the diffusion
process), `sample` (directly predicting the noisy sample`) or `v_prediction` (see section 2.4
+1 -3
View File
@@ -167,9 +167,7 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -173,9 +173,7 @@ class DDPMParallelScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2. of [Common Diffusion Noise Schedules and Sample
Steps are Flawed](https://arxiv.org/abs/2305.08891) for more information.
steps_offset (`int`, default `0`):
an offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False`, to make the last step use step 0 for the previous alpha product, as done in
stable diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -115,9 +115,7 @@ class DEISMultistepScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -178,9 +178,7 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -899,6 +897,7 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
timestep: int,
sample: torch.FloatTensor,
generator=None,
variance_noise: Optional[torch.FloatTensor] = None,
return_dict: bool = True,
) -> Union[SchedulerOutput, Tuple]:
"""
@@ -914,6 +913,9 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
A current instance of a sample created by the diffusion process.
generator (`torch.Generator`, *optional*):
A random number generator.
variance_noise (`torch.FloatTensor`):
Alternative to generating noise with `generator` by directly providing the noise for the variance
itself. Useful for methods such as [`LEdits++`].
return_dict (`bool`):
Whether or not to return a [`~schedulers.scheduling_utils.SchedulerOutput`] or `tuple`.
@@ -948,11 +950,12 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
# Upcast to avoid precision issues when computing prev_sample
sample = sample.to(torch.float32)
if self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"]:
if self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"] and variance_noise is None:
noise = randn_tensor(
model_output.shape, generator=generator, device=model_output.device, dtype=torch.float32
)
elif self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"]:
noise = variance_noise.to(device=model_output.device, dtype=torch.float32)
else:
noise = None
@@ -134,9 +134,7 @@ class DPMSolverMultistepInverseScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -792,6 +790,7 @@ class DPMSolverMultistepInverseScheduler(SchedulerMixin, ConfigMixin):
timestep: int,
sample: torch.FloatTensor,
generator=None,
variance_noise: Optional[torch.FloatTensor] = None,
return_dict: bool = True,
) -> Union[SchedulerOutput, Tuple]:
"""
@@ -807,6 +806,9 @@ class DPMSolverMultistepInverseScheduler(SchedulerMixin, ConfigMixin):
A current instance of a sample created by the diffusion process.
generator (`torch.Generator`, *optional*):
A random number generator.
variance_noise (`torch.FloatTensor`):
Alternative to generating noise with `generator` by directly providing the noise for the variance
itself. Useful for methods such as [`CycleDiffusion`].
return_dict (`bool`):
Whether or not to return a [`~schedulers.scheduling_utils.SchedulerOutput`] or `tuple`.
@@ -837,10 +839,12 @@ class DPMSolverMultistepInverseScheduler(SchedulerMixin, ConfigMixin):
self.model_outputs[i] = self.model_outputs[i + 1]
self.model_outputs[-1] = model_output
if self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"]:
if self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"] and variance_noise is None:
noise = randn_tensor(
model_output.shape, generator=generator, device=model_output.device, dtype=model_output.dtype
)
elif self.config.algorithm_type in ["sde-dpmsolver", "sde-dpmsolver++"]:
noise = variance_noise
else:
noise = None
@@ -153,9 +153,7 @@ class DPMSolverSDEScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -223,6 +223,8 @@ class DPMSolverSinglestepScheduler(SchedulerMixin, ConfigMixin):
"""
steps = num_inference_steps
order = self.config.solver_order
if order > 3:
raise ValueError("Order > 3 is not supported by this scheduler")
if self.config.lower_order_final:
if order == 3:
if steps % 3 == 0:
@@ -156,9 +156,7 @@ class EulerAncestralDiscreteScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -162,9 +162,7 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -101,9 +101,7 @@ class HeunDiscreteScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -99,9 +99,7 @@ class KDPM2AncestralDiscreteScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -98,9 +98,7 @@ class KDPM2DiscreteScheduler(SchedulerMixin, ConfigMixin):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
steps_offset (`int`, defaults to 0):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
An offset added to the inference steps, as required by some model families.
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]

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