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DN6 7fd1a8205b update 2025-08-14 14:03:01 +05:30
Sayak Paul 09e063c145 Merge branch 'main' into local-model-info 2025-08-13 21:19:54 +05:30
sayakpaul 2a9734f014 empty 2025-08-13 20:46:04 +05:30
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sayakpaul 04cd2dc451 reviewer feedback. 2025-08-13 14:50:50 +05:30
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Sayak Paul 01784c39cb Merge branch 'main' into local-model-info 2025-08-13 14:16:43 +05:30
Sayak Paul 832de66a8d Merge branch 'main' into local-model-info 2025-08-13 08:02:21 +05:30
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Sayak Paul 71843a0c8b Merge branch 'main' into local-model-info 2025-08-12 20:20:33 +05:30
Sayak Paul d1174740bb Merge branch 'main' into local-model-info 2025-08-07 10:08:33 +05:30
Sayak Paul 85279dfeee Merge branch 'main' into local-model-info 2025-08-01 08:13:57 +05:30
Sayak Paul 2d993b71d5 Merge branch 'main' into local-model-info 2025-07-29 13:58:33 +05:30
sayakpaul f38a64443f Revert "tighten compilation tests for quantization"
This reverts commit 8d431dc967.
2025-07-28 20:19:38 +05:30
sayakpaul d5c1772dc3 up 2025-07-28 20:17:24 +05:30
sayakpaul 69920eff3e feat: model_info but local. 2025-07-28 15:16:53 +05:30
sayakpaul 8d431dc967 tighten compilation tests for quantization 2025-07-28 13:27:20 +05:30
32 changed files with 671 additions and 1669 deletions
+2 -2
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@@ -5,9 +5,9 @@
- local: installation
title: Installation
- local: quicktour
title: Quickstart
title: Quicktour
- local: stable_diffusion
title: Basic performance
title: Effective and efficient diffusion
- title: DiffusionPipeline
isExpanded: false
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Qwen-Image from the Qwen team is an image generation foundation model in the Qwen series that achieves significant advances in complex text rendering and precise image editing. Experiments show strong general capabilities in both image generation and editing, with exceptional performance in text rendering, especially for Chinese.
Qwen-Image comes in the following variants:
| model type | model id |
|:----------:|:--------:|
| Qwen-Image | [`Qwen/Qwen-Image`](https://huggingface.co/Qwen/Qwen-Image) |
| Qwen-Image-Edit | [`Qwen/Qwen-Image-Edit`](https://huggingface.co/Qwen/Qwen-Image-Edit) |
Check out the model card [here](https://huggingface.co/Qwen/Qwen-Image) to learn more.
<Tip>
@@ -92,6 +87,10 @@ image.save("qwen_fewsteps.png")
- all
- __call__
## QwenImagePipelineOutput
[[autodoc]] pipelines.qwenimage.pipeline_output.QwenImagePipelineOutput
## QwenImageImg2ImgPipeline
[[autodoc]] QwenImageImg2ImgPipeline
@@ -103,13 +102,3 @@ image.save("qwen_fewsteps.png")
[[autodoc]] QwenImageInpaintPipeline
- all
- __call__
## QwenImageEditPipeline
[[autodoc]] QwenImageEditPipeline
- all
- __call__
## QwenImagePipelineOutput
[[autodoc]] pipelines.qwenimage.pipeline_output.QwenImagePipelineOutput
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@@ -10,220 +10,314 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Quickstart
[[open-in-colab]]
Diffusers is a library for developers and researchers that provides an easy inference API for generating images, videos and audio, as well as the building blocks for implementing new workflows.
# Quicktour
Diffusers provides many optimizations out-of-the-box that makes it possible to load and run large models on setups with limited memory or to accelerate inference.
Diffusion models are trained to denoise random Gaussian noise step-by-step to generate a sample of interest, such as an image or audio. This has sparked a tremendous amount of interest in generative AI, and you have probably seen examples of diffusion generated images on the internet. 🧨 Diffusers is a library aimed at making diffusion models widely accessible to everyone.
This Quickstart will give you an overview of Diffusers and get you up and generating quickly.
Whether you're a developer or an everyday user, this quicktour will introduce you to 🧨 Diffusers and help you get up and generating quickly! There are three main components of the library to know about:
> [!TIP]
> Before you begin, make sure you have a Hugging Face [account](https://huggingface.co/join) in order to use gated models like [Flux](https://huggingface.co/black-forest-labs/FLUX.1-dev).
* The [`DiffusionPipeline`] is a high-level end-to-end class designed to rapidly generate samples from pretrained diffusion models for inference.
* Popular pretrained [model](./api/models) architectures and modules that can be used as building blocks for creating diffusion systems.
* Many different [schedulers](./api/schedulers/overview) - algorithms that control how noise is added for training, and how to generate denoised images during inference.
Follow the [Installation](./installation) guide to install Diffusers if it's not already installed.
The quicktour will show you how to use the [`DiffusionPipeline`] for inference, and then walk you through how to combine a model and scheduler to replicate what's happening inside the [`DiffusionPipeline`].
<Tip>
The quicktour is a simplified version of the introductory 🧨 Diffusers [notebook](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/diffusers_intro.ipynb) to help you get started quickly. If you want to learn more about 🧨 Diffusers' goal, design philosophy, and additional details about its core API, check out the notebook!
</Tip>
Before you begin, make sure you have all the necessary libraries installed:
```py
# uncomment to install the necessary libraries in Colab
#!pip install --upgrade diffusers accelerate transformers
```
- [🤗 Accelerate](https://huggingface.co/docs/accelerate/index) speeds up model loading for inference and training.
- [🤗 Transformers](https://huggingface.co/docs/transformers/index) is required to run the most popular diffusion models, such as [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview).
## DiffusionPipeline
A diffusion model combines multiple components to generate outputs in any modality based on an input, such as a text description, image or both.
The [`DiffusionPipeline`] is the easiest way to use a pretrained diffusion system for inference. It is an end-to-end system containing the model and the scheduler. You can use the [`DiffusionPipeline`] out-of-the-box for many tasks. Take a look at the table below for some supported tasks, and for a complete list of supported tasks, check out the [🧨 Diffusers Summary](./api/pipelines/overview#diffusers-summary) table.
For a standard text-to-image model:
| **Task** | **Description** | **Pipeline**
|------------------------------|--------------------------------------------------------------------------------------------------------------|-----------------|
| Unconditional Image Generation | generate an image from Gaussian noise | [unconditional_image_generation](./using-diffusers/unconditional_image_generation) |
| Text-Guided Image Generation | generate an image given a text prompt | [conditional_image_generation](./using-diffusers/conditional_image_generation) |
| Text-Guided Image-to-Image Translation | adapt an image guided by a text prompt | [img2img](./using-diffusers/img2img) |
| Text-Guided Image-Inpainting | fill the masked part of an image given the image, the mask and a text prompt | [inpaint](./using-diffusers/inpaint) |
| Text-Guided Depth-to-Image Translation | adapt parts of an image guided by a text prompt while preserving structure via depth estimation | [depth2img](./using-diffusers/depth2img) |
1. A text encoder turns a prompt into embeddings that guide the denoising process. Some models have more than one text encoder.
2. A scheduler contains the algorithmic specifics for gradually denoising initial random noise into clean outputs. Different schedulers affect generation speed and quality.
3. A UNet or diffusion transformer (DiT) is the workhorse of a diffusion model.
Start by creating an instance of a [`DiffusionPipeline`] and specify which pipeline checkpoint you would like to download.
You can use the [`DiffusionPipeline`] for any [checkpoint](https://huggingface.co/models?library=diffusers&sort=downloads) stored on the Hugging Face Hub.
In this quicktour, you'll load the [`stable-diffusion-v1-5`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) checkpoint for text-to-image generation.
At each step, it performs the denoising predictions, such as how much noise to remove or the general direction in which to steer the noise to generate better quality outputs.
<Tip warning={true}>
The UNet or DiT repeats this loop for a set amount of steps to generate the final output.
4. A variational autoencoder (VAE) encodes and decodes pixels to a spatially compressed latent-space. *Latents* are compressed representations of an image and are more efficient to work with. The UNet or DiT operates on latents, and the clean latents at the end are decoded back into images.
For [Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion) models, please carefully read the [license](https://huggingface.co/spaces/CompVis/stable-diffusion-license) first before running the model. 🧨 Diffusers implements a [`safety_checker`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/safety_checker.py) to prevent offensive or harmful content, but the model's improved image generation capabilities can still produce potentially harmful content.
The [`DiffusionPipeline`] packages all these components into a single class for inference. There are several arguments in [`~DiffusionPipeline.__call__`] you can change, such as `num_inference_steps`, that affect the diffusion process. Try different values and arguments to see how they change generation quality or speed.
</Tip>
Load a model with [`~DiffusionPipeline.from_pretrained`] and describe what you'd like to generate. The example below uses the default argument values.
Load the model with the [`~DiffusionPipeline.from_pretrained`] method:
<hfoptions id="diffusionpipeline">
<hfoption id="text-to-image">
```python
>>> from diffusers import DiffusionPipeline
Use `.images[0]` to access the generated image output.
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"Qwen/Qwen-Image", torch_dtype=torch.bfloat16, device_map="cuda"
)
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
>>> pipeline = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", use_safetensors=True)
```
</hfoption>
<hfoption id="text-to-video">
Use `.frames[0]` to access the generated video output and [`~utils.export_to_video`] to save the video.
The [`DiffusionPipeline`] downloads and caches all modeling, tokenization, and scheduling components. You'll see that the Stable Diffusion pipeline is composed of the [`UNet2DConditionModel`] and [`PNDMScheduler`] among other things:
```py
import torch
from diffusers import AutoencoderKLWan, DiffusionPipeline
from diffusers.quantizers import PipelineQuantizationConfig
from diffusers.utils import export_to_video
vae = AutoencoderKLWan.from_pretrained(
"Wan-AI/Wan2.2-T2V-A14B-Diffusers",
subfolder="vae",
torch_dtype=torch.float32
)
pipeline = DiffusionPipeline.from_pretrained(
"Wan-AI/Wan2.2-T2V-A14B-Diffusers",
vae=vae
torch_dtype=torch.bfloat16,
device_map="cuda"
)
prompt = """
Cinematic video of a sleek cat lounging on a colorful inflatable in a crystal-clear turquoise pool in Palm Springs,
sipping a salt-rimmed margarita through a straw. Golden-hour sunlight glows over mid-century modern homes and swaying palms.
Shot in rich Sony a7S III: with moody, glamorous color grading, subtle lens flares, and soft vintage film grain.
Ripples shimmer as a warm desert breeze stirs the water, blending luxury and playful charm in an epic, gorgeously composed frame.
"""
video = pipeline(prompt=prompt, num_frames=81, num_inference_steps=40).frames[0]
export_to_video(video, "output.mp4", fps=16)
>>> pipeline
StableDiffusionPipeline {
"_class_name": "StableDiffusionPipeline",
"_diffusers_version": "0.21.4",
...,
"scheduler": [
"diffusers",
"PNDMScheduler"
],
...,
"unet": [
"diffusers",
"UNet2DConditionModel"
],
"vae": [
"diffusers",
"AutoencoderKL"
]
}
```
</hfoption>
</hfoptions>
We strongly recommend running the pipeline on a GPU because the model consists of roughly 1.4 billion parameters.
You can move the generator object to a GPU, just like you would in PyTorch:
## LoRA
```python
>>> pipeline.to("cuda")
```
Adapters insert a small number of trainable parameters to the original base model. Only the inserted parameters are fine-tuned while the rest of the model weights remain frozen. This makes it fast and cheap to fine-tune a model on a new style. Among adapters, [LoRA's](./tutorials/using_peft_for_inference) are the most popular.
Now you can pass a text prompt to the `pipeline` to generate an image, and then access the denoised image. By default, the image output is wrapped in a [`PIL.Image`](https://pillow.readthedocs.io/en/stable/reference/Image.html?highlight=image#the-image-class) object.
Add a LoRA to a pipeline with the [`~loaders.QwenImageLoraLoaderMixin.load_lora_weights`] method. Some LoRA's require a special word to trigger it, such as `Realism`, in the example below. Check a LoRA's model card to see if it requires a trigger word.
```python
>>> image = pipeline("An image of a squirrel in Picasso style").images[0]
>>> image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/image_of_squirrel_painting.png"/>
</div>
Save the image by calling `save`:
```python
>>> image.save("image_of_squirrel_painting.png")
```
### Local pipeline
You can also use the pipeline locally. The only difference is you need to download the weights first:
```bash
!git lfs install
!git clone https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5
```
Then load the saved weights into the pipeline:
```python
>>> pipeline = DiffusionPipeline.from_pretrained("./stable-diffusion-v1-5", use_safetensors=True)
```
Now, you can run the pipeline as you would in the section above.
### Swapping schedulers
Different schedulers come with different denoising speeds and quality trade-offs. The best way to find out which one works best for you is to try them out! One of the main features of 🧨 Diffusers is to allow you to easily switch between schedulers. For example, to replace the default [`PNDMScheduler`] with the [`EulerDiscreteScheduler`], load it with the [`~diffusers.ConfigMixin.from_config`] method:
```py
import torch
from diffusers import DiffusionPipeline
>>> from diffusers import EulerDiscreteScheduler
pipeline = DiffusionPipeline.from_pretrained(
"Qwen/Qwen-Image", torch_dtype=torch.bfloat16, device_map="cuda"
)
pipeline.load_lora_weights(
"flymy-ai/qwen-image-realism-lora",
)
prompt = """
super Realism cinematic film still of a cat sipping a margarita in a pool in Palm Springs in the style of umempart, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
>>> pipeline = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", use_safetensors=True)
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
```
Check out the [LoRA](./tutorials/using_peft_for_inference) docs or Adapters section to learn more.
Try generating an image with the new scheduler and see if you notice a difference!
## Quantization
In the next section, you'll take a closer look at the components - the model and scheduler - that make up the [`DiffusionPipeline`] and learn how to use these components to generate an image of a cat.
[Quantization](./quantization/overview) stores data in fewer bits to reduce memory usage. It may also speed up inference because it takes less time to perform calculations with fewer bits.
## Models
Diffusers provides several quantization backends and picking one depends on your use case. For example, [bitsandbytes](./quantization/bitsandbytes) and [torchao](./quantization/torchao) are both simple and easy to use for inference, but torchao supports more [quantization types](./quantization/torchao#supported-quantization-types) like fp8.
Most models take a noisy sample, and at each timestep it predicts the *noise residual* (other models learn to predict the previous sample directly or the velocity or [`v-prediction`](https://github.com/huggingface/diffusers/blob/5e5ce13e2f89ac45a0066cb3f369462a3cf1d9ef/src/diffusers/schedulers/scheduling_ddim.py#L110)), the difference between a less noisy image and the input image. You can mix and match models to create other diffusion systems.
Configure [`PipelineQuantizationConfig`] with the backend to use, the specific arguments (refer to the [API](./api/quantization) reference for available arguments) for that backend, and which components to quantize. The example below quantizes the model to 4-bits and only uses 14.93GB of memory.
Models are initiated with the [`~ModelMixin.from_pretrained`] method which also locally caches the model weights so it is faster the next time you load the model. For the quicktour, you'll load the [`UNet2DModel`], a basic unconditional image generation model with a checkpoint trained on cat images:
```py
import torch
from diffusers import DiffusionPipeline
from diffusers.quantizers import PipelineQuantizationConfig
>>> from diffusers import UNet2DModel
quant_config = PipelineQuantizationConfig(
quant_backend="bitsandbytes_4bit",
quant_kwargs={"load_in_4bit": True, "bnb_4bit_quant_type": "nf4", "bnb_4bit_compute_dtype": torch.bfloat16},
components_to_quantize=["transformer", "text_encoder"],
)
pipeline = DiffusionPipeline.from_pretrained(
"Qwen/Qwen-Image",
torch_dtype=torch.bfloat16,
quantization_config=quant_config,
device_map="cuda"
)
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
>>> repo_id = "google/ddpm-cat-256"
>>> model = UNet2DModel.from_pretrained(repo_id, use_safetensors=True)
```
Take a look at the [Quantization](./quantization/overview) section for more details.
> [!TIP]
> Use the [`AutoModel`] API to automatically select a model class if you're unsure of which one to use.
## Optimizations
Modern diffusion models are very large and have billions of parameters. The iterative denoising process is also computationally intensive and slow. Diffusers provides techniques for reducing memory usage and boosting inference speed. These techniques can be combined with quantization to optimize for both memory usage and inference speed.
### Memory usage
The text encoders and UNet or DiT can use up as much as ~30GB of memory, exceeding the amount available on many free-tier or consumer GPUs.
Offloading stores weights that aren't currently used on the CPU and only moves them to the GPU when they're needed. There are a few offloading types and the example below uses [model offloading](./optimization/memory#model-offloading). This moves an entire model, like a text encoder or transformer, to the CPU when it isn't actively being used.
Call [`~DiffusionPipeline.enable_model_cpu_offload`] to activate it. By combining quantization and offloading, the following example only requires ~12.54GB of memory.
To access the model parameters, call `model.config`:
```py
import torch
from diffusers import DiffusionPipeline
from diffusers.quantizers import PipelineQuantizationConfig
quant_config = PipelineQuantizationConfig(
quant_backend="bitsandbytes_4bit",
quant_kwargs={"load_in_4bit": True, "bnb_4bit_quant_type": "nf4", "bnb_4bit_compute_dtype": torch.bfloat16},
components_to_quantize=["transformer", "text_encoder"],
)
pipeline = DiffusionPipeline.from_pretrained(
"Qwen/Qwen-Image",
torch_dtype=torch.bfloat16,
quantization_config=quant_config,
device_map="cuda"
)
pipeline.enable_model_cpu_offload()
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
>>> model.config
```
Refer to the [Reduce memory usage](./optimization/memory) docs to learn more about other memory reducing techniques.
The model configuration is a 🧊 frozen 🧊 dictionary, which means those parameters can't be changed after the model is created. This is intentional and ensures that the parameters used to define the model architecture at the start remain the same, while other parameters can still be adjusted during inference.
### Inference speed
Some of the most important parameters are:
The denoising loop performs a lot of computations and can be slow. Methods like [torch.compile](./optimization/fp16#torchcompile) increases inference speed by compiling the computations into an optimized kernel. Compilation is slow for the first generation but successive generations should be much faster.
* `sample_size`: the height and width dimension of the input sample.
* `in_channels`: the number of input channels of the input sample.
* `down_block_types` and `up_block_types`: the type of down- and upsampling blocks used to create the UNet architecture.
* `block_out_channels`: the number of output channels of the downsampling blocks; also used in reverse order for the number of input channels of the upsampling blocks.
* `layers_per_block`: the number of ResNet blocks present in each UNet block.
The example below uses [regional compilation](./optimization/fp16#regional-compilation) to only compile small regions of a model. It reduces cold-start latency while also providing a runtime speed up.
Call [`~ModelMixin.compile_repeated_blocks`] on the model to activate it.
To use the model for inference, create the image shape with random Gaussian noise. It should have a `batch` axis because the model can receive multiple random noises, a `channel` axis corresponding to the number of input channels, and a `sample_size` axis for the height and width of the image:
```py
import torch
from diffusers import DiffusionPipeline
>>> import torch
pipeline = DiffusionPipeline.from_pretrained(
"Qwen/Qwen-Image", torch_dtype=torch.bfloat16, device_map="cuda"
)
>>> torch.manual_seed(0)
pipeline.transformer.compile_repeated_blocks(
fullgraph=True,
)
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
>>> noisy_sample = torch.randn(1, model.config.in_channels, model.config.sample_size, model.config.sample_size)
>>> noisy_sample.shape
torch.Size([1, 3, 256, 256])
```
Check out the [Accelerate inference](./optimization/fp16) or [Caching](./optimization/cache) docs for more methods that speed up inference.
For inference, pass the noisy image and a `timestep` to the model. The `timestep` indicates how noisy the input image is, with more noise at the beginning and less at the end. This helps the model determine its position in the diffusion process, whether it is closer to the start or the end. Use the `sample` method to get the model output:
```py
>>> with torch.no_grad():
... noisy_residual = model(sample=noisy_sample, timestep=2).sample
```
To generate actual examples though, you'll need a scheduler to guide the denoising process. In the next section, you'll learn how to couple a model with a scheduler.
## Schedulers
Schedulers manage going from a noisy sample to a less noisy sample given the model output - in this case, it is the `noisy_residual`.
<Tip>
🧨 Diffusers is a toolbox for building diffusion systems. While the [`DiffusionPipeline`] is a convenient way to get started with a pre-built diffusion system, you can also choose your own model and scheduler components separately to build a custom diffusion system.
</Tip>
For the quicktour, you'll instantiate the [`DDPMScheduler`] with its [`~diffusers.ConfigMixin.from_config`] method:
```py
>>> from diffusers import DDPMScheduler
>>> scheduler = DDPMScheduler.from_pretrained(repo_id)
>>> scheduler
DDPMScheduler {
"_class_name": "DDPMScheduler",
"_diffusers_version": "0.21.4",
"beta_end": 0.02,
"beta_schedule": "linear",
"beta_start": 0.0001,
"clip_sample": true,
"clip_sample_range": 1.0,
"dynamic_thresholding_ratio": 0.995,
"num_train_timesteps": 1000,
"prediction_type": "epsilon",
"sample_max_value": 1.0,
"steps_offset": 0,
"thresholding": false,
"timestep_spacing": "leading",
"trained_betas": null,
"variance_type": "fixed_small"
}
```
<Tip>
💡 Unlike a model, a scheduler does not have trainable weights and is parameter-free!
</Tip>
Some of the most important parameters are:
* `num_train_timesteps`: the length of the denoising process or, in other words, the number of timesteps required to process random Gaussian noise into a data sample.
* `beta_schedule`: the type of noise schedule to use for inference and training.
* `beta_start` and `beta_end`: the start and end noise values for the noise schedule.
To predict a slightly less noisy image, pass the following to the scheduler's [`~diffusers.DDPMScheduler.step`] method: model output, `timestep`, and current `sample`.
```py
>>> less_noisy_sample = scheduler.step(model_output=noisy_residual, timestep=2, sample=noisy_sample).prev_sample
>>> less_noisy_sample.shape
torch.Size([1, 3, 256, 256])
```
The `less_noisy_sample` can be passed to the next `timestep` where it'll get even less noisy! Let's bring it all together now and visualize the entire denoising process.
First, create a function that postprocesses and displays the denoised image as a `PIL.Image`:
```py
>>> import PIL.Image
>>> import numpy as np
>>> def display_sample(sample, i):
... image_processed = sample.cpu().permute(0, 2, 3, 1)
... image_processed = (image_processed + 1.0) * 127.5
... image_processed = image_processed.numpy().astype(np.uint8)
... image_pil = PIL.Image.fromarray(image_processed[0])
... display(f"Image at step {i}")
... display(image_pil)
```
To speed up the denoising process, move the input and model to a GPU:
```py
>>> model.to("cuda")
>>> noisy_sample = noisy_sample.to("cuda")
```
Now create a denoising loop that predicts the residual of the less noisy sample, and computes the less noisy sample with the scheduler:
```py
>>> import tqdm
>>> sample = noisy_sample
>>> for i, t in enumerate(tqdm.tqdm(scheduler.timesteps)):
... # 1. predict noise residual
... with torch.no_grad():
... residual = model(sample, t).sample
... # 2. compute less noisy image and set x_t -> x_t-1
... sample = scheduler.step(residual, t, sample).prev_sample
... # 3. optionally look at image
... if (i + 1) % 50 == 0:
... display_sample(sample, i + 1)
```
Sit back and watch as a cat is generated from nothing but noise! 😻
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/diffusion-quicktour.png"/>
</div>
## Next steps
Hopefully, you generated some cool images with 🧨 Diffusers in this quicktour! For your next steps, you can:
* Train or finetune a model to generate your own images in the [training](./tutorials/basic_training) tutorial.
* See example official and community [training or finetuning scripts](https://github.com/huggingface/diffusers/tree/main/examples#-diffusers-examples) for a variety of use cases.
* Learn more about loading, accessing, changing, and comparing schedulers in the [Using different Schedulers](./using-diffusers/schedulers) guide.
* Explore prompt engineering, speed and memory optimizations, and tips and tricks for generating higher-quality images with the [Stable Diffusion](./stable_diffusion) guide.
* Dive deeper into speeding up 🧨 Diffusers with guides on [optimized PyTorch on a GPU](./optimization/fp16), and inference guides for running [Stable Diffusion on Apple Silicon (M1/M2)](./optimization/mps) and [ONNX Runtime](./optimization/onnx).
+215 -86
View File
@@ -10,123 +10,252 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Effective and efficient diffusion
[[open-in-colab]]
# Basic performance
Getting the [`DiffusionPipeline`] to generate images in a certain style or include what you want can be tricky. Often times, you have to run the [`DiffusionPipeline`] several times before you end up with an image you're happy with. But generating something out of nothing is a computationally intensive process, especially if you're running inference over and over again.
Diffusion is a random process that is computationally demanding. You may need to run the [`DiffusionPipeline`] several times before getting a desired output. That's why it's important to carefully balance generation speed and memory usage in order to iterate faster,
This is why it's important to get the most *computational* (speed) and *memory* (GPU vRAM) efficiency from the pipeline to reduce the time between inference cycles so you can iterate faster.
This guide recommends some basic performance tips for using the [`DiffusionPipeline`]. Refer to the Inference Optimization section docs such as [Accelerate inference](./optimization/fp16) or [Reduce memory usage](./optimization/memory) for more detailed performance guides.
This tutorial walks you through how to generate faster and better with the [`DiffusionPipeline`].
## Memory usage
Begin by loading the [`stable-diffusion-v1-5/stable-diffusion-v1-5`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) model:
Reducing the amount of memory used indirectly speeds up generation and can help a model fit on device.
The [`~DiffusionPipeline.enable_model_cpu_offload`] method moves a model to the CPU when it is not in use to save GPU memory.
```py
import torch
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.bfloat16,
device_map="cuda"
)
pipeline.enable_model_cpu_offload()
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
pipeline(prompt).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
pipeline = DiffusionPipeline.from_pretrained(model_id, use_safetensors=True)
```
## Inference speed
The example prompt you'll use is a portrait of an old warrior chief, but feel free to use your own prompt:
Denoising is the most computationally demanding process during diffusion. Methods that optimizes this process accelerates inference speed. Try the following methods for a speed up.
```python
prompt = "portrait photo of a old warrior chief"
```
- Add `device_map="cuda"` to place the pipeline on a GPU. Placing a model on an accelerator, like a GPU, increases speed because it performs computations in parallel.
- Set `torch_dtype=torch.bfloat16` to execute the pipeline in half-precision. Reducing the data type precision increases speed because it takes less time to perform computations in a lower precision.
## Speed
```py
<Tip>
💡 If you don't have access to a GPU, you can use one for free from a GPU provider like [Colab](https://colab.research.google.com/)!
</Tip>
One of the simplest ways to speed up inference is to place the pipeline on a GPU the same way you would with any PyTorch module:
```python
pipeline = pipeline.to("cuda")
```
To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reusing_seeds):
```python
import torch
import time
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.bfloat16,
device_map="cuda
)
generator = torch.Generator("cuda").manual_seed(0)
```
- Use a faster scheduler, such as [`DPMSolverMultistepScheduler`], which only requires ~20-25 steps.
- Set `num_inference_steps` to a lower value. Reducing the number of inference steps reduces the overall number of computations. However, this can result in lower generation quality.
Now you can generate an image:
```python
image = pipeline(prompt, generator=generator).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_1.png">
</div>
This process took ~30 seconds on a T4 GPU (it might be faster if your allocated GPU is better than a T4). By default, the [`DiffusionPipeline`] runs inference with full `float32` precision for 50 inference steps. You can speed this up by switching to a lower precision like `float16` or running fewer inference steps.
Let's start by loading the model in `float16` and generate an image:
```python
import torch
pipeline = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16, use_safetensors=True)
pipeline = pipeline.to("cuda")
generator = torch.Generator("cuda").manual_seed(0)
image = pipeline(prompt, generator=generator).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_2.png">
</div>
This time, it only took ~11 seconds to generate the image, which is almost 3x faster than before!
<Tip>
💡 We strongly suggest always running your pipelines in `float16`, and so far, we've rarely seen any degradation in output quality.
</Tip>
Another option is to reduce the number of inference steps. Choosing a more efficient scheduler could help decrease the number of steps without sacrificing output quality. You can find which schedulers are compatible with the current model in the [`DiffusionPipeline`] by calling the `compatibles` method:
```python
pipeline.scheduler.compatibles
[
diffusers.schedulers.scheduling_lms_discrete.LMSDiscreteScheduler,
diffusers.schedulers.scheduling_unipc_multistep.UniPCMultistepScheduler,
diffusers.schedulers.scheduling_k_dpm_2_discrete.KDPM2DiscreteScheduler,
diffusers.schedulers.scheduling_deis_multistep.DEISMultistepScheduler,
diffusers.schedulers.scheduling_euler_discrete.EulerDiscreteScheduler,
diffusers.schedulers.scheduling_dpmsolver_multistep.DPMSolverMultistepScheduler,
diffusers.schedulers.scheduling_ddpm.DDPMScheduler,
diffusers.schedulers.scheduling_dpmsolver_singlestep.DPMSolverSinglestepScheduler,
diffusers.schedulers.scheduling_k_dpm_2_ancestral_discrete.KDPM2AncestralDiscreteScheduler,
diffusers.utils.dummy_torch_and_torchsde_objects.DPMSolverSDEScheduler,
diffusers.schedulers.scheduling_heun_discrete.HeunDiscreteScheduler,
diffusers.schedulers.scheduling_pndm.PNDMScheduler,
diffusers.schedulers.scheduling_euler_ancestral_discrete.EulerAncestralDiscreteScheduler,
diffusers.schedulers.scheduling_ddim.DDIMScheduler,
]
```
The Stable Diffusion model uses the [`PNDMScheduler`] by default which usually requires ~50 inference steps, but more performant schedulers like [`DPMSolverMultistepScheduler`], require only ~20 or 25 inference steps. Use the [`~ConfigMixin.from_config`] method to load a new scheduler:
```python
from diffusers import DPMSolverMultistepScheduler
```py
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
start_time = time.perf_counter()
image = pipeline(prompt).images[0]
end_time = time.perf_counter()
print(f"Image generation took {end_time - start_time:.3f} seconds")
```
## Generation quality
Now set the `num_inference_steps` to 20:
Many modern diffusion models deliver high-quality images out-of-the-box. However, you can still improve generation quality by trying the following.
```python
generator = torch.Generator("cuda").manual_seed(0)
image = pipeline(prompt, generator=generator, num_inference_steps=20).images[0]
image
```
- Try a more detailed and descriptive prompt. Include details such as the image medium, subject, style, and aesthetic. A negative prompt may also help by guiding a model away from undesirable features by using words like low quality or blurry.
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_3.png">
</div>
```py
import torch
from diffusers import DiffusionPipeline
Great, you've managed to cut the inference time to just 4 seconds! ⚡️
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.bfloat16,
device_map="cuda"
)
## Memory
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
negative_prompt = "low quality, blurry, ugly, poor details"
pipeline(prompt, negative_prompt=negative_prompt).images[0]
```
The other key to improving pipeline performance is consuming less memory, which indirectly implies more speed, since you're often trying to maximize the number of images generated per second. The easiest way to see how many images you can generate at once is to try out different batch sizes until you get an `OutOfMemoryError` (OOM).
For more details about creating better prompts, take a look at the [Prompt techniques](./using-diffusers/weighted_prompts) doc.
Create a function that'll generate a batch of images from a list of prompts and `Generators`. Make sure to assign each `Generator` a seed so you can reuse it if it produces a good result.
- Try a different scheduler, like [`HeunDiscreteScheduler`] or [`LMSDiscreteScheduler`], that gives up generation speed for quality.
```python
def get_inputs(batch_size=1):
generator = [torch.Generator("cuda").manual_seed(i) for i in range(batch_size)]
prompts = batch_size * [prompt]
num_inference_steps = 20
```py
import torch
from diffusers import DiffusionPipeline, HeunDiscreteScheduler
return {"prompt": prompts, "generator": generator, "num_inference_steps": num_inference_steps}
```
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.bfloat16,
device_map="cuda"
)
pipeline.scheduler = HeunDiscreteScheduler.from_config(pipeline.scheduler.config)
Start with `batch_size=4` and see how much memory you've consumed:
prompt = """
cinematic film still of a cat sipping a margarita in a pool in Palm Springs, California
highly detailed, high budget hollywood movie, cinemascope, moody, epic, gorgeous, film grain
"""
negative_prompt = "low quality, blurry, ugly, poor details"
pipeline(prompt, negative_prompt=negative_prompt).images[0]
```
```python
from diffusers.utils import make_image_grid
images = pipeline(**get_inputs(batch_size=4)).images
make_image_grid(images, 2, 2)
```
Unless you have a GPU with more vRAM, the code above probably returned an `OOM` error! Most of the memory is taken up by the cross-attention layers. Instead of running this operation in a batch, you can run it sequentially to save a significant amount of memory. All you have to do is configure the pipeline to use the [`~DiffusionPipeline.enable_attention_slicing`] function:
```python
pipeline.enable_attention_slicing()
```
Now try increasing the `batch_size` to 8!
```python
images = pipeline(**get_inputs(batch_size=8)).images
make_image_grid(images, rows=2, cols=4)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_5.png">
</div>
Whereas before you couldn't even generate a batch of 4 images, now you can generate a batch of 8 images at ~3.5 seconds per image! This is probably the fastest you can go on a T4 GPU without sacrificing quality.
## Quality
In the last two sections, you learned how to optimize the speed of your pipeline by using `fp16`, reducing the number of inference steps by using a more performant scheduler, and enabling attention slicing to reduce memory consumption. Now you're going to focus on how to improve the quality of generated images.
### Better checkpoints
The most obvious step is to use better checkpoints. The Stable Diffusion model is a good starting point, and since its official launch, several improved versions have also been released. However, using a newer version doesn't automatically mean you'll get better results. You'll still have to experiment with different checkpoints yourself, and do a little research (such as using [negative prompts](https://minimaxir.com/2022/11/stable-diffusion-negative-prompt/)) to get the best results.
As the field grows, there are more and more high-quality checkpoints finetuned to produce certain styles. Try exploring the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) and [Diffusers Gallery](https://huggingface.co/spaces/huggingface-projects/diffusers-gallery) to find one you're interested in!
### Better pipeline components
You can also try replacing the current pipeline components with a newer version. Let's try loading the latest [autoencoder](https://huggingface.co/stabilityai/stable-diffusion-2-1/tree/main/vae) from Stability AI into the pipeline, and generate some images:
```python
from diffusers import AutoencoderKL
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16).to("cuda")
pipeline.vae = vae
images = pipeline(**get_inputs(batch_size=8)).images
make_image_grid(images, rows=2, cols=4)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_6.png">
</div>
### Better prompt engineering
The text prompt you use to generate an image is super important, so much so that it is called *prompt engineering*. Some considerations to keep during prompt engineering are:
- How is the image or similar images of the one I want to generate stored on the internet?
- What additional detail can I give that steers the model towards the style I want?
With this in mind, let's improve the prompt to include color and higher quality details:
```python
prompt += ", tribal panther make up, blue on red, side profile, looking away, serious eyes"
prompt += " 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta"
```
Generate a batch of images with the new prompt:
```python
images = pipeline(**get_inputs(batch_size=8)).images
make_image_grid(images, rows=2, cols=4)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_7.png">
</div>
Pretty impressive! Let's tweak the second image - corresponding to the `Generator` with a seed of `1` - a bit more by adding some text about the age of the subject:
```python
prompts = [
"portrait photo of the oldest warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of an old warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of a warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of a young warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
]
generator = [torch.Generator("cuda").manual_seed(1) for _ in range(len(prompts))]
images = pipeline(prompt=prompts, generator=generator, num_inference_steps=25).images
make_image_grid(images, 2, 2)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_8.png">
</div>
## Next steps
Diffusers offers more advanced and powerful optimizations such as [group-offloading](./optimization/memory#group-offloading) and [regional compilation](./optimization/fp16#regional-compilation). To learn more about how to maximize performance, take a look at the Inference Optimization section.
In this tutorial, you learned how to optimize a [`DiffusionPipeline`] for computational and memory efficiency as well as improving the quality of generated outputs. If you're interested in making your pipeline even faster, take a look at the following resources:
- Learn how [PyTorch 2.0](./optimization/fp16) and [`torch.compile`](https://pytorch.org/docs/stable/generated/torch.compile.html) can yield 5 - 300% faster inference speed. On an A100 GPU, inference can be up to 50% faster!
- If you can't use PyTorch 2, we recommend you install [xFormers](./optimization/xformers). Its memory-efficient attention mechanism works great with PyTorch 1.13.1 for faster speed and reduced memory consumption.
- Other optimization techniques, such as model offloading, are covered in [this guide](./optimization/fp16).
-24
View File
@@ -112,30 +112,6 @@ print(pipe.transformer.dtype, pipe.vae.dtype) # (torch.bfloat16, torch.float16)
If a component is not explicitly specified in the dictionary and no `default` is provided, it will be loaded with `torch.float32`.
### Parallel loading
Large models are often [sharded](../training/distributed_inference#model-sharding) into smaller files so that they are easier to load. Diffusers supports loading shards in parallel to speed up the loading process.
Set the environment variables below to enable parallel loading.
- Set `HF_ENABLE_PARALLEL_LOADING` to `"YES"` to enable parallel loading of shards.
- Set `HF_PARALLEL_LOADING_WORKERS` to configure the number of parallel threads to use when loading shards. More workers loads a model faster but uses more memory.
The `device_map` argument should be set to `"cuda"` to pre-allocate a large chunk of memory based on the model size. This substantially reduces model load time because warming up the memory allocator now avoids many smaller calls to the allocator later.
```py
import os
import torch
from diffusers import DiffusionPipeline
os.environ["HF_ENABLE_PARALLEL_LOADING"] = "YES"
pipeline = DiffusionPipeline.from_pretrained(
"Wan-AI/Wan2.2-I2V-A14B-Diffusers",
torch_dtype=torch.bfloat16,
device_map="cuda"
)
```
### Local pipeline
To load a pipeline locally, use [git-lfs](https://git-lfs.github.com/) to manually download a checkpoint to your local disk.
-2
View File
@@ -489,7 +489,6 @@ else:
"PixArtAlphaPipeline",
"PixArtSigmaPAGPipeline",
"PixArtSigmaPipeline",
"QwenImageEditPipeline",
"QwenImageImg2ImgPipeline",
"QwenImageInpaintPipeline",
"QwenImagePipeline",
@@ -1124,7 +1123,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
PixArtAlphaPipeline,
PixArtSigmaPAGPipeline,
PixArtSigmaPipeline,
QwenImageEditPipeline,
QwenImageImg2ImgPipeline,
QwenImageInpaintPipeline,
QwenImagePipeline,
+6 -11
View File
@@ -17,6 +17,7 @@
import functools
import importlib
import inspect
import math
import os
from array import array
from collections import OrderedDict, defaultdict
@@ -716,33 +717,27 @@ def _expand_device_map(device_map, param_names):
# Adapted from: https://github.com/huggingface/transformers/blob/0687d481e2c71544501ef9cb3eef795a6e79b1de/src/transformers/modeling_utils.py#L5859
def _caching_allocator_warmup(
model, expanded_device_map: Dict[str, torch.device], dtype: torch.dtype, hf_quantizer: Optional[DiffusersQuantizer]
) -> None:
def _caching_allocator_warmup(model, expanded_device_map: Dict[str, torch.device], dtype: torch.dtype) -> None:
"""
This function warm-ups the caching allocator based on the size of the model tensors that will reside on each
device. It allows to have one large call to Malloc, instead of recursively calling it later when loading the model,
which is actually the loading speed bottleneck. Calling this function allows to cut the model loading time by a
very large margin.
"""
factor = 2 if hf_quantizer is None else hf_quantizer.get_cuda_warm_up_factor()
# Remove disk and cpu devices, and cast to proper torch.device
accelerator_device_map = {
param: torch.device(device)
for param, device in expanded_device_map.items()
if str(device) not in ["cpu", "disk"]
}
total_byte_count = defaultdict(lambda: 0)
parameter_count = defaultdict(lambda: 0)
for param_name, device in accelerator_device_map.items():
try:
param = model.get_parameter(param_name)
except AttributeError:
param = model.get_buffer(param_name)
# The dtype of different parameters may be different with composite models or `keep_in_fp32_modules`
param_byte_count = param.numel() * param.element_size()
# TODO: account for TP when needed.
total_byte_count[device] += param_byte_count
parameter_count[device] += math.prod(param.shape)
# This will kick off the caching allocator to avoid having to Malloc afterwards
for device, byte_count in total_byte_count.items():
_ = torch.empty(byte_count // factor, dtype=dtype, device=device, requires_grad=False)
for device, param_count in parameter_count.items():
_ = torch.empty(param_count, dtype=dtype, device=device, requires_grad=False)
+6 -4
View File
@@ -42,8 +42,9 @@ from ..quantizers import DiffusersAutoQuantizer, DiffusersQuantizer
from ..quantizers.quantization_config import QuantizationMethod
from ..utils import (
CONFIG_NAME,
ENV_VARS_TRUE_VALUES,
FLAX_WEIGHTS_NAME,
HF_ENABLE_PARALLEL_LOADING,
HF_PARALLEL_LOADING_FLAG,
SAFE_WEIGHTS_INDEX_NAME,
SAFETENSORS_WEIGHTS_NAME,
WEIGHTS_INDEX_NAME,
@@ -961,7 +962,7 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
dduf_entries: Optional[Dict[str, DDUFEntry]] = kwargs.pop("dduf_entries", None)
disable_mmap = kwargs.pop("disable_mmap", False)
is_parallel_loading_enabled = HF_ENABLE_PARALLEL_LOADING
is_parallel_loading_enabled = os.environ.get(HF_PARALLEL_LOADING_FLAG, "").upper() in ENV_VARS_TRUE_VALUES
if is_parallel_loading_enabled and not low_cpu_mem_usage:
raise NotImplementedError("Parallel loading is not supported when not using `low_cpu_mem_usage`.")
@@ -1532,9 +1533,10 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
# tensors using their expected shape and not performing any initialization of the memory (empty data).
# When the actual device allocations happen, the allocator already has a pool of unused device memory
# that it can re-use for faster loading of the model.
if device_map is not None:
# TODO: add support for warmup with hf_quantizer
if device_map is not None and hf_quantizer is None:
expanded_device_map = _expand_device_map(device_map, expected_keys)
_caching_allocator_warmup(model, expanded_device_map, dtype, hf_quantizer)
_caching_allocator_warmup(model, expanded_device_map, dtype)
offload_index = {} if device_map is not None and "disk" in device_map.values() else None
state_dict_folder, state_dict_index = None, None
@@ -28,7 +28,7 @@ from ..cache_utils import CacheMixin
from ..embeddings import CogView3CombinedTimestepSizeEmbeddings
from ..modeling_outputs import Transformer2DModelOutput
from ..modeling_utils import ModelMixin
from ..normalization import LayerNorm, RMSNorm
from ..normalization import AdaLayerNormContinuous
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
@@ -584,38 +584,6 @@ class CogView4RotaryPosEmbed(nn.Module):
return (freqs.cos(), freqs.sin())
class CogView4AdaLayerNormContinuous(nn.Module):
"""
CogView4-only final AdaLN: LN(x) -> Linear(cond) -> chunk -> affine. Matches Megatron: **no activation** before the
Linear on conditioning embedding.
"""
def __init__(
self,
embedding_dim: int,
conditioning_embedding_dim: int,
elementwise_affine: bool = True,
eps: float = 1e-5,
bias: bool = True,
norm_type: str = "layer_norm",
):
super().__init__()
self.linear = nn.Linear(conditioning_embedding_dim, embedding_dim * 2, bias=bias)
if norm_type == "layer_norm":
self.norm = LayerNorm(embedding_dim, eps, elementwise_affine, bias)
elif norm_type == "rms_norm":
self.norm = RMSNorm(embedding_dim, eps, elementwise_affine)
else:
raise ValueError(f"unknown norm_type {norm_type}")
def forward(self, x: torch.Tensor, conditioning_embedding: torch.Tensor) -> torch.Tensor:
# *** NO SiLU here ***
emb = self.linear(conditioning_embedding.to(x.dtype))
scale, shift = torch.chunk(emb, 2, dim=1)
x = self.norm(x) * (1 + scale)[:, None, :] + shift[:, None, :]
return x
class CogView4Transformer2DModel(ModelMixin, ConfigMixin, PeftAdapterMixin, CacheMixin):
r"""
Args:
@@ -698,7 +666,7 @@ class CogView4Transformer2DModel(ModelMixin, ConfigMixin, PeftAdapterMixin, Cach
)
# 4. Output projection
self.norm_out = CogView4AdaLayerNormContinuous(inner_dim, time_embed_dim, elementwise_affine=False)
self.norm_out = AdaLayerNormContinuous(inner_dim, time_embed_dim, elementwise_affine=False)
self.proj_out = nn.Linear(inner_dim, patch_size * patch_size * out_channels, bias=True)
self.gradient_checkpointing = False
@@ -12,6 +12,7 @@
# See the License for the specific language governing permissions and
# limitations under the License.
import functools
import math
from typing import Any, Dict, List, Optional, Tuple, Union
@@ -160,9 +161,9 @@ class QwenEmbedRope(nn.Module):
super().__init__()
self.theta = theta
self.axes_dim = axes_dim
pos_index = torch.arange(4096)
neg_index = torch.arange(4096).flip(0) * -1 - 1
self.pos_freqs = torch.cat(
pos_index = torch.arange(1024)
neg_index = torch.arange(1024).flip(0) * -1 - 1
pos_freqs = torch.cat(
[
self.rope_params(pos_index, self.axes_dim[0], self.theta),
self.rope_params(pos_index, self.axes_dim[1], self.theta),
@@ -170,7 +171,7 @@ class QwenEmbedRope(nn.Module):
],
dim=1,
)
self.neg_freqs = torch.cat(
neg_freqs = torch.cat(
[
self.rope_params(neg_index, self.axes_dim[0], self.theta),
self.rope_params(neg_index, self.axes_dim[1], self.theta),
@@ -179,8 +180,10 @@ class QwenEmbedRope(nn.Module):
dim=1,
)
self.rope_cache = {}
self.register_buffer("pos_freqs", pos_freqs, persistent=False)
self.register_buffer("neg_freqs", neg_freqs, persistent=False)
# DO NOT USING REGISTER BUFFER HERE, IT WILL CAUSE COMPLEX NUMBERS LOSE ITS IMAGINARY PART
# 是否使用 scale rope
self.scale_rope = scale_rope
def rope_params(self, index, dim, theta=10000):
@@ -198,48 +201,35 @@ class QwenEmbedRope(nn.Module):
Args: video_fhw: [frame, height, width] a list of 3 integers representing the shape of the video Args:
txt_length: [bs] a list of 1 integers representing the length of the text
"""
if self.pos_freqs.device != device:
self.pos_freqs = self.pos_freqs.to(device)
self.neg_freqs = self.neg_freqs.to(device)
if isinstance(video_fhw, list):
video_fhw = video_fhw[0]
if not isinstance(video_fhw, list):
video_fhw = [video_fhw]
frame, height, width = video_fhw
rope_key = f"{frame}_{height}_{width}"
vid_freqs = []
max_vid_index = 0
for idx, fhw in enumerate(video_fhw):
frame, height, width = fhw
rope_key = f"{idx}_{height}_{width}"
if not torch.compiler.is_compiling():
if rope_key not in self.rope_cache:
self.rope_cache[rope_key] = self._compute_video_freqs(frame, height, width)
vid_freqs = self.rope_cache[rope_key]
else:
vid_freqs = self._compute_video_freqs(frame, height, width)
if not torch.compiler.is_compiling():
if rope_key not in self.rope_cache:
self.rope_cache[rope_key] = self._compute_video_freqs(frame, height, width, idx)
video_freq = self.rope_cache[rope_key]
else:
video_freq = self._compute_video_freqs(frame, height, width, idx)
video_freq = video_freq.to(device)
vid_freqs.append(video_freq)
if self.scale_rope:
max_vid_index = max(height // 2, width // 2, max_vid_index)
else:
max_vid_index = max(height, width, max_vid_index)
if self.scale_rope:
max_vid_index = max(height // 2, width // 2)
else:
max_vid_index = max(height, width)
max_len = max(txt_seq_lens)
txt_freqs = self.pos_freqs[max_vid_index : max_vid_index + max_len, ...]
vid_freqs = torch.cat(vid_freqs, dim=0)
return vid_freqs, txt_freqs
@functools.lru_cache(maxsize=None)
def _compute_video_freqs(self, frame, height, width, idx=0):
def _compute_video_freqs(self, frame, height, width):
seq_lens = frame * height * width
freqs_pos = self.pos_freqs.split([x // 2 for x in self.axes_dim], dim=1)
freqs_neg = self.neg_freqs.split([x // 2 for x in self.axes_dim], dim=1)
freqs_frame = freqs_pos[0][idx : idx + frame].view(frame, 1, 1, -1).expand(frame, height, width, -1)
freqs_frame = freqs_pos[0][:frame].view(frame, 1, 1, -1).expand(frame, height, width, -1)
if self.scale_rope:
freqs_height = torch.cat([freqs_neg[1][-(height - height // 2) :], freqs_pos[1][: height // 2]], dim=0)
freqs_height = freqs_height.view(1, height, 1, -1).expand(frame, height, width, -1)
@@ -290,7 +290,7 @@ class ModularPipelineBlocks(ConfigMixin, PushToHubMixin):
def from_pretrained(
cls,
pretrained_model_name_or_path: str,
trust_remote_code: bool = False,
trust_remote_code: Optional[bool] = None,
**kwargs,
):
hub_kwargs_names = [
+1 -7
View File
@@ -391,7 +391,6 @@ else:
"QwenImagePipeline",
"QwenImageImg2ImgPipeline",
"QwenImageInpaintPipeline",
"QwenImageEditPipeline",
]
try:
if not is_onnx_available():
@@ -709,12 +708,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
from .paint_by_example import PaintByExamplePipeline
from .pia import PIAPipeline
from .pixart_alpha import PixArtAlphaPipeline, PixArtSigmaPipeline
from .qwenimage import (
QwenImageEditPipeline,
QwenImageImg2ImgPipeline,
QwenImageInpaintPipeline,
QwenImagePipeline,
)
from .qwenimage import QwenImageImg2ImgPipeline, QwenImageInpaintPipeline, QwenImagePipeline
from .sana import SanaControlNetPipeline, SanaPipeline, SanaSprintImg2ImgPipeline, SanaSprintPipeline
from .semantic_stable_diffusion import SemanticStableDiffusionPipeline
from .shap_e import ShapEImg2ImgPipeline, ShapEPipeline
@@ -613,9 +613,6 @@ def _assign_components_to_devices(
def _get_final_device_map(device_map, pipeline_class, passed_class_obj, init_dict, library, max_memory, **kwargs):
# TODO: seperate out different device_map methods when it gets to it.
if device_map != "balanced":
return device_map
# To avoid circular import problem.
from diffusers import pipelines
+7 -10
View File
@@ -108,7 +108,7 @@ LIBRARIES = []
for library in LOADABLE_CLASSES:
LIBRARIES.append(library)
SUPPORTED_DEVICE_MAP = ["balanced"] + [get_device()]
SUPPORTED_DEVICE_MAP = ["balanced"]
logger = logging.get_logger(__name__)
@@ -988,15 +988,12 @@ class DiffusionPipeline(ConfigMixin, PushToHubMixin):
_maybe_warn_for_wrong_component_in_quant_config(init_dict, quantization_config)
for name, (library_name, class_name) in logging.tqdm(init_dict.items(), desc="Loading pipeline components..."):
# 7.1 device_map shenanigans
if final_device_map is not None:
if isinstance(final_device_map, dict) and len(final_device_map) > 0:
component_device = final_device_map.get(name, None)
if component_device is not None:
current_device_map = {"": component_device}
else:
current_device_map = None
elif isinstance(final_device_map, str):
current_device_map = final_device_map
if final_device_map is not None and len(final_device_map) > 0:
component_device = final_device_map.get(name, None)
if component_device is not None:
current_device_map = {"": component_device}
else:
current_device_map = None
# 7.2 - now that JAX/Flax is an official framework of the library, we might load from Flax names
class_name = class_name[4:] if class_name.startswith("Flax") else class_name
@@ -24,7 +24,6 @@ except OptionalDependencyNotAvailable:
else:
_import_structure["modeling_qwenimage"] = ["ReduxImageEncoder"]
_import_structure["pipeline_qwenimage"] = ["QwenImagePipeline"]
_import_structure["pipeline_qwenimage_edit"] = ["QwenImageEditPipeline"]
_import_structure["pipeline_qwenimage_img2img"] = ["QwenImageImg2ImgPipeline"]
_import_structure["pipeline_qwenimage_inpaint"] = ["QwenImageInpaintPipeline"]
@@ -36,7 +35,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
from ...utils.dummy_torch_and_transformers_objects import * # noqa F403
else:
from .pipeline_qwenimage import QwenImagePipeline
from .pipeline_qwenimage_edit import QwenImageEditPipeline
from .pipeline_qwenimage_img2img import QwenImageImg2ImgPipeline
from .pipeline_qwenimage_inpaint import QwenImageInpaintPipeline
else:
@@ -253,9 +253,6 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if prompt_embeds is None:
prompt_embeds, prompt_embeds_mask = self._get_qwen_prompt_embeds(prompt, device)
prompt_embeds = prompt_embeds[:, :max_sequence_length]
prompt_embeds_mask = prompt_embeds_mask[:, :max_sequence_length]
_, seq_len, _ = prompt_embeds.shape
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
@@ -319,6 +316,20 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if max_sequence_length is not None and max_sequence_length > 1024:
raise ValueError(f"`max_sequence_length` cannot be greater than 1024 but is {max_sequence_length}")
@staticmethod
def _prepare_latent_image_ids(batch_size, height, width, device, dtype):
latent_image_ids = torch.zeros(height, width, 3)
latent_image_ids[..., 1] = latent_image_ids[..., 1] + torch.arange(height)[:, None]
latent_image_ids[..., 2] = latent_image_ids[..., 2] + torch.arange(width)[None, :]
latent_image_id_height, latent_image_id_width, latent_image_id_channels = latent_image_ids.shape
latent_image_ids = latent_image_ids.reshape(
latent_image_id_height * latent_image_id_width, latent_image_id_channels
)
return latent_image_ids.to(device=device, dtype=dtype)
@staticmethod
def _pack_latents(latents, batch_size, num_channels_latents, height, width):
latents = latents.view(batch_size, num_channels_latents, height // 2, 2, width // 2, 2)
@@ -391,7 +402,8 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
shape = (batch_size, 1, num_channels_latents, height, width)
if latents is not None:
return latents.to(device=device, dtype=dtype)
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents.to(device=device, dtype=dtype), latent_image_ids
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
@@ -402,7 +414,9 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
latents = self._pack_latents(latents, batch_size, num_channels_latents, height, width)
return latents
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents, latent_image_ids
@property
def guidance_scale(self):
@@ -580,7 +594,7 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
# 4. Prepare latent variables
num_channels_latents = self.transformer.config.in_channels // 4
latents = self.prepare_latents(
latents, latent_image_ids = self.prepare_latents(
batch_size * num_images_per_prompt,
num_channels_latents,
height,
@@ -590,7 +604,7 @@ class QwenImagePipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
generator,
latents,
)
img_shapes = [[(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)]] * batch_size
img_shapes = [(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)] * batch_size
# 5. Prepare timesteps
sigmas = np.linspace(1.0, 1 / num_inference_steps, num_inference_steps) if sigmas is None else sigmas
@@ -1,877 +0,0 @@
# Copyright 2025 Qwen-Image Team and The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import inspect
import math
from typing import Any, Callable, Dict, List, Optional, Union
import numpy as np
import torch
from transformers import Qwen2_5_VLForConditionalGeneration, Qwen2Tokenizer, Qwen2VLProcessor
from ...image_processor import PipelineImageInput, VaeImageProcessor
from ...loaders import QwenImageLoraLoaderMixin
from ...models import AutoencoderKLQwenImage, QwenImageTransformer2DModel
from ...schedulers import FlowMatchEulerDiscreteScheduler
from ...utils import is_torch_xla_available, logging, replace_example_docstring
from ...utils.torch_utils import randn_tensor
from ..pipeline_utils import DiffusionPipeline
from .pipeline_output import QwenImagePipelineOutput
if is_torch_xla_available():
import torch_xla.core.xla_model as xm
XLA_AVAILABLE = True
else:
XLA_AVAILABLE = False
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
EXAMPLE_DOC_STRING = """
Examples:
```py
>>> import torch
>>> from PIL import Image
>>> from diffusers import QwenImageEditPipeline
>>> from diffusers.utils import load_image
>>> pipe = QwenImageEditPipeline.from_pretrained("Qwen/Qwen-Image-Edit", torch_dtype=torch.bfloat16)
>>> pipe.to("cuda")
>>> image = load_image(
... "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/yarn-art-pikachu.png"
... ).convert("RGB")
>>> prompt = (
... "Make Pikachu hold a sign that says 'Qwen Edit is awesome', yarn art style, detailed, vibrant colors"
... )
>>> # Depending on the variant being used, the pipeline call will slightly vary.
>>> # Refer to the pipeline documentation for more details.
>>> image = pipe(image, prompt, num_inference_steps=50).images[0]
>>> image.save("qwenimage_edit.png")
```
"""
PREFERRED_QWENIMAGE_RESOLUTIONS = [
(672, 1568),
(688, 1504),
(720, 1456),
(752, 1392),
(800, 1328),
(832, 1248),
(880, 1184),
(944, 1104),
(1024, 1024),
(1104, 944),
(1184, 880),
(1248, 832),
(1328, 800),
(1392, 752),
(1456, 720),
(1504, 688),
(1568, 672),
]
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.calculate_shift
def calculate_shift(
image_seq_len,
base_seq_len: int = 256,
max_seq_len: int = 4096,
base_shift: float = 0.5,
max_shift: float = 1.15,
):
m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
b = base_shift - m * base_seq_len
mu = image_seq_len * m + b
return mu
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.retrieve_timesteps
def retrieve_timesteps(
scheduler,
num_inference_steps: Optional[int] = None,
device: Optional[Union[str, torch.device]] = None,
timesteps: Optional[List[int]] = None,
sigmas: Optional[List[float]] = None,
**kwargs,
):
r"""
Calls the scheduler's `set_timesteps` method and retrieves timesteps from the scheduler after the call. Handles
custom timesteps. Any kwargs will be supplied to `scheduler.set_timesteps`.
Args:
scheduler (`SchedulerMixin`):
The scheduler to get timesteps from.
num_inference_steps (`int`):
The number of diffusion steps used when generating samples with a pre-trained model. If used, `timesteps`
must be `None`.
device (`str` or `torch.device`, *optional*):
The device to which the timesteps should be moved to. If `None`, the timesteps are not moved.
timesteps (`List[int]`, *optional*):
Custom timesteps used to override the timestep spacing strategy of the scheduler. If `timesteps` is passed,
`num_inference_steps` and `sigmas` must be `None`.
sigmas (`List[float]`, *optional*):
Custom sigmas used to override the timestep spacing strategy of the scheduler. If `sigmas` is passed,
`num_inference_steps` and `timesteps` must be `None`.
Returns:
`Tuple[torch.Tensor, int]`: A tuple where the first element is the timestep schedule from the scheduler and the
second element is the number of inference steps.
"""
if timesteps is not None and sigmas is not None:
raise ValueError("Only one of `timesteps` or `sigmas` can be passed. Please choose one to set custom values")
if timesteps is not None:
accepts_timesteps = "timesteps" in set(inspect.signature(scheduler.set_timesteps).parameters.keys())
if not accepts_timesteps:
raise ValueError(
f"The current scheduler class {scheduler.__class__}'s `set_timesteps` does not support custom"
f" timestep schedules. Please check whether you are using the correct scheduler."
)
scheduler.set_timesteps(timesteps=timesteps, device=device, **kwargs)
timesteps = scheduler.timesteps
num_inference_steps = len(timesteps)
elif sigmas is not None:
accept_sigmas = "sigmas" in set(inspect.signature(scheduler.set_timesteps).parameters.keys())
if not accept_sigmas:
raise ValueError(
f"The current scheduler class {scheduler.__class__}'s `set_timesteps` does not support custom"
f" sigmas schedules. Please check whether you are using the correct scheduler."
)
scheduler.set_timesteps(sigmas=sigmas, device=device, **kwargs)
timesteps = scheduler.timesteps
num_inference_steps = len(timesteps)
else:
scheduler.set_timesteps(num_inference_steps, device=device, **kwargs)
timesteps = scheduler.timesteps
return timesteps, num_inference_steps
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_img2img.retrieve_latents
def retrieve_latents(
encoder_output: torch.Tensor, generator: Optional[torch.Generator] = None, sample_mode: str = "sample"
):
if hasattr(encoder_output, "latent_dist") and sample_mode == "sample":
return encoder_output.latent_dist.sample(generator)
elif hasattr(encoder_output, "latent_dist") and sample_mode == "argmax":
return encoder_output.latent_dist.mode()
elif hasattr(encoder_output, "latents"):
return encoder_output.latents
else:
raise AttributeError("Could not access latents of provided encoder_output")
def calculate_dimensions(target_area, ratio):
width = math.sqrt(target_area * ratio)
height = width / ratio
width = round(width / 32) * 32
height = round(height / 32) * 32
return width, height, None
class QwenImageEditPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
r"""
The Qwen-Image-Edit pipeline for image editing.
Args:
transformer ([`QwenImageTransformer2DModel`]):
Conditional Transformer (MMDiT) architecture to denoise the encoded image latents.
scheduler ([`FlowMatchEulerDiscreteScheduler`]):
A scheduler to be used in combination with `transformer` to denoise the encoded image latents.
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
text_encoder ([`Qwen2.5-VL-7B-Instruct`]):
[Qwen2.5-VL-7B-Instruct](https://huggingface.co/Qwen/Qwen2.5-VL-7B-Instruct), specifically the
[Qwen2.5-VL-7B-Instruct](https://huggingface.co/Qwen/Qwen2.5-VL-7B-Instruct) variant.
tokenizer (`QwenTokenizer`):
Tokenizer of class
[CLIPTokenizer](https://huggingface.co/docs/transformers/en/model_doc/clip#transformers.CLIPTokenizer).
"""
model_cpu_offload_seq = "text_encoder->transformer->vae"
_callback_tensor_inputs = ["latents", "prompt_embeds"]
def __init__(
self,
scheduler: FlowMatchEulerDiscreteScheduler,
vae: AutoencoderKLQwenImage,
text_encoder: Qwen2_5_VLForConditionalGeneration,
tokenizer: Qwen2Tokenizer,
processor: Qwen2VLProcessor,
transformer: QwenImageTransformer2DModel,
):
super().__init__()
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
processor=processor,
transformer=transformer,
scheduler=scheduler,
)
self.vae_scale_factor = 2 ** len(self.vae.temperal_downsample) if getattr(self, "vae", None) else 8
self.latent_channels = self.vae.config.z_dim if getattr(self, "vae", None) else 16
# QwenImage latents are turned into 2x2 patches and packed. This means the latent width and height has to be divisible
# by the patch size. So the vae scale factor is multiplied by the patch size to account for this
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor * 2)
self.vl_processor = processor
self.tokenizer_max_length = 1024
self.prompt_template_encode = "<|im_start|>system\nDescribe the key features of the input image (color, shape, size, texture, objects, background), then explain how the user's text instruction should alter or modify the image. Generate a new image that meets the user's requirements while maintaining consistency with the original input where appropriate.<|im_end|>\n<|im_start|>user\n<|vision_start|><|image_pad|><|vision_end|>{}<|im_end|>\n<|im_start|>assistant\n"
self.prompt_template_encode_start_idx = 64
self.default_sample_size = 128
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._extract_masked_hidden
def _extract_masked_hidden(self, hidden_states: torch.Tensor, mask: torch.Tensor):
bool_mask = mask.bool()
valid_lengths = bool_mask.sum(dim=1)
selected = hidden_states[bool_mask]
split_result = torch.split(selected, valid_lengths.tolist(), dim=0)
return split_result
def _get_qwen_prompt_embeds(
self,
prompt: Union[str, List[str]] = None,
image: Optional[torch.Tensor] = None,
device: Optional[torch.device] = None,
dtype: Optional[torch.dtype] = None,
):
device = device or self._execution_device
dtype = dtype or self.text_encoder.dtype
prompt = [prompt] if isinstance(prompt, str) else prompt
template = self.prompt_template_encode
drop_idx = self.prompt_template_encode_start_idx
txt = [template.format(e) for e in prompt]
model_inputs = self.processor(
text=txt,
images=image,
padding=True,
return_tensors="pt",
).to(device)
outputs = self.text_encoder(
input_ids=model_inputs.input_ids,
attention_mask=model_inputs.attention_mask,
pixel_values=model_inputs.pixel_values,
image_grid_thw=model_inputs.image_grid_thw,
output_hidden_states=True,
)
hidden_states = outputs.hidden_states[-1]
split_hidden_states = self._extract_masked_hidden(hidden_states, model_inputs.attention_mask)
split_hidden_states = [e[drop_idx:] for e in split_hidden_states]
attn_mask_list = [torch.ones(e.size(0), dtype=torch.long, device=e.device) for e in split_hidden_states]
max_seq_len = max([e.size(0) for e in split_hidden_states])
prompt_embeds = torch.stack(
[torch.cat([u, u.new_zeros(max_seq_len - u.size(0), u.size(1))]) for u in split_hidden_states]
)
encoder_attention_mask = torch.stack(
[torch.cat([u, u.new_zeros(max_seq_len - u.size(0))]) for u in attn_mask_list]
)
prompt_embeds = prompt_embeds.to(dtype=dtype, device=device)
return prompt_embeds, encoder_attention_mask
def encode_prompt(
self,
prompt: Union[str, List[str]],
image: Optional[torch.Tensor] = None,
device: Optional[torch.device] = None,
num_images_per_prompt: int = 1,
prompt_embeds: Optional[torch.Tensor] = None,
prompt_embeds_mask: Optional[torch.Tensor] = None,
max_sequence_length: int = 1024,
):
r"""
Args:
prompt (`str` or `List[str]`, *optional*):
prompt to be encoded
image (`torch.Tensor`, *optional*):
image to be encoded
device: (`torch.device`):
torch device
num_images_per_prompt (`int`):
number of images that should be generated per prompt
prompt_embeds (`torch.Tensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not
provided, text embeddings will be generated from `prompt` input argument.
"""
device = device or self._execution_device
prompt = [prompt] if isinstance(prompt, str) else prompt
batch_size = len(prompt) if prompt_embeds is None else prompt_embeds.shape[0]
if prompt_embeds is None:
prompt_embeds, prompt_embeds_mask = self._get_qwen_prompt_embeds(prompt, image, device)
_, seq_len, _ = prompt_embeds.shape
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
prompt_embeds_mask = prompt_embeds_mask.repeat(1, num_images_per_prompt, 1)
prompt_embeds_mask = prompt_embeds_mask.view(batch_size * num_images_per_prompt, seq_len)
return prompt_embeds, prompt_embeds_mask
def check_inputs(
self,
prompt,
height,
width,
negative_prompt=None,
prompt_embeds=None,
negative_prompt_embeds=None,
prompt_embeds_mask=None,
negative_prompt_embeds_mask=None,
callback_on_step_end_tensor_inputs=None,
max_sequence_length=None,
):
if height % (self.vae_scale_factor * 2) != 0 or width % (self.vae_scale_factor * 2) != 0:
logger.warning(
f"`height` and `width` have to be divisible by {self.vae_scale_factor * 2} but are {height} and {width}. Dimensions will be resized accordingly"
)
if callback_on_step_end_tensor_inputs is not None and not all(
k in self._callback_tensor_inputs for k in callback_on_step_end_tensor_inputs
):
raise ValueError(
f"`callback_on_step_end_tensor_inputs` has to be in {self._callback_tensor_inputs}, but found {[k for k in callback_on_step_end_tensor_inputs if k not in self._callback_tensor_inputs]}"
)
if prompt is not None and prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to"
" only forward one of the two."
)
elif prompt is None and prompt_embeds is None:
raise ValueError(
"Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined."
)
elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if negative_prompt is not None and negative_prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:"
f" {negative_prompt_embeds}. Please make sure to only forward one of the two."
)
if prompt_embeds is not None and prompt_embeds_mask is None:
raise ValueError(
"If `prompt_embeds` are provided, `prompt_embeds_mask` also have to be passed. Make sure to generate `prompt_embeds_mask` from the same text encoder that was used to generate `prompt_embeds`."
)
if negative_prompt_embeds is not None and negative_prompt_embeds_mask is None:
raise ValueError(
"If `negative_prompt_embeds` are provided, `negative_prompt_embeds_mask` also have to be passed. Make sure to generate `negative_prompt_embeds_mask` from the same text encoder that was used to generate `negative_prompt_embeds`."
)
if max_sequence_length is not None and max_sequence_length > 1024:
raise ValueError(f"`max_sequence_length` cannot be greater than 1024 but is {max_sequence_length}")
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._pack_latents
def _pack_latents(latents, batch_size, num_channels_latents, height, width):
latents = latents.view(batch_size, num_channels_latents, height // 2, 2, width // 2, 2)
latents = latents.permute(0, 2, 4, 1, 3, 5)
latents = latents.reshape(batch_size, (height // 2) * (width // 2), num_channels_latents * 4)
return latents
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._unpack_latents
def _unpack_latents(latents, height, width, vae_scale_factor):
batch_size, num_patches, channels = latents.shape
# VAE applies 8x compression on images but we must also account for packing which requires
# latent height and width to be divisible by 2.
height = 2 * (int(height) // (vae_scale_factor * 2))
width = 2 * (int(width) // (vae_scale_factor * 2))
latents = latents.view(batch_size, height // 2, width // 2, channels // 4, 2, 2)
latents = latents.permute(0, 3, 1, 4, 2, 5)
latents = latents.reshape(batch_size, channels // (2 * 2), 1, height, width)
return latents
def _encode_vae_image(self, image: torch.Tensor, generator: torch.Generator):
if isinstance(generator, list):
image_latents = [
retrieve_latents(self.vae.encode(image[i : i + 1]), generator=generator[i], sample_mode="argmax")
for i in range(image.shape[0])
]
image_latents = torch.cat(image_latents, dim=0)
else:
image_latents = retrieve_latents(self.vae.encode(image), generator=generator, sample_mode="argmax")
latents_mean = (
torch.tensor(self.vae.config.latents_mean)
.view(1, self.latent_channels, 1, 1, 1)
.to(image_latents.device, image_latents.dtype)
)
latents_std = (
torch.tensor(self.vae.config.latents_std)
.view(1, self.latent_channels, 1, 1, 1)
.to(image_latents.device, image_latents.dtype)
)
image_latents = (image_latents - latents_mean) / latents_std
return image_latents
def enable_vae_slicing(self):
r"""
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
"""
self.vae.enable_slicing()
def disable_vae_slicing(self):
r"""
Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_slicing()
def enable_vae_tiling(self):
r"""
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow
processing larger images.
"""
self.vae.enable_tiling()
def disable_vae_tiling(self):
r"""
Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_tiling()
def prepare_latents(
self,
image,
batch_size,
num_channels_latents,
height,
width,
dtype,
device,
generator,
latents=None,
):
# VAE applies 8x compression on images but we must also account for packing which requires
# latent height and width to be divisible by 2.
height = 2 * (int(height) // (self.vae_scale_factor * 2))
width = 2 * (int(width) // (self.vae_scale_factor * 2))
shape = (batch_size, 1, num_channels_latents, height, width)
image_latents = None
if image is not None:
image = image.to(device=device, dtype=dtype)
if image.shape[1] != self.latent_channels:
image_latents = self._encode_vae_image(image=image, generator=generator)
else:
image_latents = image
if batch_size > image_latents.shape[0] and batch_size % image_latents.shape[0] == 0:
# expand init_latents for batch_size
additional_image_per_prompt = batch_size // image_latents.shape[0]
image_latents = torch.cat([image_latents] * additional_image_per_prompt, dim=0)
elif batch_size > image_latents.shape[0] and batch_size % image_latents.shape[0] != 0:
raise ValueError(
f"Cannot duplicate `image` of batch size {image_latents.shape[0]} to {batch_size} text prompts."
)
else:
image_latents = torch.cat([image_latents], dim=0)
image_latent_height, image_latent_width = image_latents.shape[3:]
image_latents = self._pack_latents(
image_latents, batch_size, num_channels_latents, image_latent_height, image_latent_width
)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
if latents is None:
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
latents = self._pack_latents(latents, batch_size, num_channels_latents, height, width)
else:
latents = latents.to(device=device, dtype=dtype)
return latents, image_latents
@property
def guidance_scale(self):
return self._guidance_scale
@property
def attention_kwargs(self):
return self._attention_kwargs
@property
def num_timesteps(self):
return self._num_timesteps
@property
def current_timestep(self):
return self._current_timestep
@property
def interrupt(self):
return self._interrupt
@torch.no_grad()
@replace_example_docstring(EXAMPLE_DOC_STRING)
def __call__(
self,
image: Optional[PipelineImageInput] = None,
prompt: Union[str, List[str]] = None,
negative_prompt: Union[str, List[str]] = None,
true_cfg_scale: float = 4.0,
height: Optional[int] = None,
width: Optional[int] = None,
num_inference_steps: int = 50,
sigmas: Optional[List[float]] = None,
guidance_scale: float = 1.0,
num_images_per_prompt: int = 1,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
latents: Optional[torch.Tensor] = None,
prompt_embeds: Optional[torch.Tensor] = None,
prompt_embeds_mask: Optional[torch.Tensor] = None,
negative_prompt_embeds: Optional[torch.Tensor] = None,
negative_prompt_embeds_mask: Optional[torch.Tensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
attention_kwargs: Optional[Dict[str, Any]] = None,
callback_on_step_end: Optional[Callable[[int, int, Dict], None]] = None,
callback_on_step_end_tensor_inputs: List[str] = ["latents"],
max_sequence_length: int = 512,
_auto_resize: bool = True,
):
r"""
Function invoked when calling the pipeline for generation.
Args:
prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`.
instead.
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. If not defined, one has to pass
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `true_cfg_scale` is
not greater than `1`).
true_cfg_scale (`float`, *optional*, defaults to 1.0):
When > 1.0 and a provided `negative_prompt`, enables true classifier-free guidance.
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
The height in pixels of the generated image. This is set to 1024 by default for the best results.
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
The width in pixels of the generated image. This is set to 1024 by default for the best results.
num_inference_steps (`int`, *optional*, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
sigmas (`List[float]`, *optional*):
Custom sigmas to use for the denoising process with schedulers which support a `sigmas` argument in
their `set_timesteps` method. If not defined, the default behavior when `num_inference_steps` is passed
will be used.
guidance_scale (`float`, *optional*, defaults to 3.5):
Guidance scale as defined in [Classifier-Free Diffusion
Guidance](https://huggingface.co/papers/2207.12598). `guidance_scale` is defined as `w` of equation 2.
of [Imagen Paper](https://huggingface.co/papers/2205.11487). Guidance scale is enabled by setting
`guidance_scale > 1`. Higher guidance scale encourages to generate images that are closely linked to
the text `prompt`, usually at the expense of lower image quality.
num_images_per_prompt (`int`, *optional*, defaults to 1):
The number of images to generate per prompt.
generator (`torch.Generator` or `List[torch.Generator]`, *optional*):
One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html)
to make generation deterministic.
latents (`torch.Tensor`, *optional*):
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will be generated by sampling using the supplied random `generator`.
prompt_embeds (`torch.Tensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not
provided, text embeddings will be generated from `prompt` input argument.
negative_prompt_embeds (`torch.Tensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input
argument.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generate image. Choose between
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~pipelines.qwenimage.QwenImagePipelineOutput`] instead of a plain tuple.
attention_kwargs (`dict`, *optional*):
A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under
`self.processor` in
[diffusers.models.attention_processor](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
callback_on_step_end (`Callable`, *optional*):
A function that calls at the end of each denoising steps during the inference. The function is called
with the following arguments: `callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int,
callback_kwargs: Dict)`. `callback_kwargs` will include a list of all tensors as specified by
`callback_on_step_end_tensor_inputs`.
callback_on_step_end_tensor_inputs (`List`, *optional*):
The list of tensor inputs for the `callback_on_step_end` function. The tensors specified in the list
will be passed as `callback_kwargs` argument. You will only be able to include variables listed in the
`._callback_tensor_inputs` attribute of your pipeline class.
max_sequence_length (`int` defaults to 512): Maximum sequence length to use with the `prompt`.
Examples:
Returns:
[`~pipelines.qwenimage.QwenImagePipelineOutput`] or `tuple`:
[`~pipelines.qwenimage.QwenImagePipelineOutput`] if `return_dict` is True, otherwise a `tuple`. When
returning a tuple, the first element is a list with the generated images.
"""
image_size = image[0].size if isinstance(image, list) else image.size
width, height = image_size
calculated_width, calculated_height, _ = calculate_dimensions(1024 * 1024, width / height)
height = height or calculated_height
width = width or calculated_width
multiple_of = self.vae_scale_factor * 2
width = width // multiple_of * multiple_of
height = height // multiple_of * multiple_of
# 1. Check inputs. Raise error if not correct
self.check_inputs(
prompt,
height,
width,
negative_prompt=negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
prompt_embeds_mask=prompt_embeds_mask,
negative_prompt_embeds_mask=negative_prompt_embeds_mask,
callback_on_step_end_tensor_inputs=callback_on_step_end_tensor_inputs,
max_sequence_length=max_sequence_length,
)
self._guidance_scale = guidance_scale
self._attention_kwargs = attention_kwargs
self._current_timestep = None
self._interrupt = False
# 2. Define call parameters
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
device = self._execution_device
# 3. Preprocess image
if image is not None and not (isinstance(image, torch.Tensor) and image.size(1) == self.latent_channels):
img = image[0] if isinstance(image, list) else image
image_height, image_width = self.image_processor.get_default_height_width(img)
aspect_ratio = image_width / image_height
if _auto_resize:
_, image_width, image_height = min(
(abs(aspect_ratio - w / h), w, h) for w, h in PREFERRED_QWENIMAGE_RESOLUTIONS
)
image_width = image_width // multiple_of * multiple_of
image_height = image_height // multiple_of * multiple_of
image = self.image_processor.resize(image, image_height, image_width)
prompt_image = image
image = self.image_processor.preprocess(image, image_height, image_width)
image = image.unsqueeze(2)
has_neg_prompt = negative_prompt is not None or (
negative_prompt_embeds is not None and negative_prompt_embeds_mask is not None
)
do_true_cfg = true_cfg_scale > 1 and has_neg_prompt
prompt_embeds, prompt_embeds_mask = self.encode_prompt(
image=prompt_image,
prompt=prompt,
prompt_embeds=prompt_embeds,
prompt_embeds_mask=prompt_embeds_mask,
device=device,
num_images_per_prompt=num_images_per_prompt,
max_sequence_length=max_sequence_length,
)
if do_true_cfg:
# negative image is the same size as the original image, but all pixels are white
# negative_image = Image.new("RGB", (image.width, image.height), (255, 255, 255))
negative_prompt_embeds, negative_prompt_embeds_mask = self.encode_prompt(
image=prompt_image,
prompt=negative_prompt,
prompt_embeds=negative_prompt_embeds,
prompt_embeds_mask=negative_prompt_embeds_mask,
device=device,
num_images_per_prompt=num_images_per_prompt,
max_sequence_length=max_sequence_length,
)
# 4. Prepare latent variables
num_channels_latents = self.transformer.config.in_channels // 4
latents, image_latents = self.prepare_latents(
image,
batch_size * num_images_per_prompt,
num_channels_latents,
height,
width,
prompt_embeds.dtype,
device,
generator,
latents,
)
img_shapes = [
[
(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2),
(1, image_height // self.vae_scale_factor // 2, image_width // self.vae_scale_factor // 2),
]
] * batch_size
# 5. Prepare timesteps
sigmas = np.linspace(1.0, 1 / num_inference_steps, num_inference_steps) if sigmas is None else sigmas
image_seq_len = latents.shape[1]
mu = calculate_shift(
image_seq_len,
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
)
timesteps, num_inference_steps = retrieve_timesteps(
self.scheduler,
num_inference_steps,
device,
sigmas=sigmas,
mu=mu,
)
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0)
self._num_timesteps = len(timesteps)
# handle guidance
if self.transformer.config.guidance_embeds:
guidance = torch.full([1], guidance_scale, device=device, dtype=torch.float32)
guidance = guidance.expand(latents.shape[0])
else:
guidance = None
if self.attention_kwargs is None:
self._attention_kwargs = {}
txt_seq_lens = prompt_embeds_mask.sum(dim=1).tolist() if prompt_embeds_mask is not None else None
negative_txt_seq_lens = (
negative_prompt_embeds_mask.sum(dim=1).tolist() if negative_prompt_embeds_mask is not None else None
)
# 6. Denoising loop
self.scheduler.set_begin_index(0)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
if self.interrupt:
continue
self._current_timestep = t
latent_model_input = latents
if image_latents is not None:
latent_model_input = torch.cat([latents, image_latents], dim=1)
# broadcast to batch dimension in a way that's compatible with ONNX/Core ML
timestep = t.expand(latents.shape[0]).to(latents.dtype)
with self.transformer.cache_context("cond"):
noise_pred = self.transformer(
hidden_states=latent_model_input,
timestep=timestep / 1000,
guidance=guidance,
encoder_hidden_states_mask=prompt_embeds_mask,
encoder_hidden_states=prompt_embeds,
img_shapes=img_shapes,
txt_seq_lens=txt_seq_lens,
attention_kwargs=self.attention_kwargs,
return_dict=False,
)[0]
noise_pred = noise_pred[:, : latents.size(1)]
if do_true_cfg:
with self.transformer.cache_context("uncond"):
neg_noise_pred = self.transformer(
hidden_states=latent_model_input,
timestep=timestep / 1000,
guidance=guidance,
encoder_hidden_states_mask=negative_prompt_embeds_mask,
encoder_hidden_states=negative_prompt_embeds,
img_shapes=img_shapes,
txt_seq_lens=negative_txt_seq_lens,
attention_kwargs=self.attention_kwargs,
return_dict=False,
)[0]
neg_noise_pred = neg_noise_pred[:, : latents.size(1)]
comb_pred = neg_noise_pred + true_cfg_scale * (noise_pred - neg_noise_pred)
cond_norm = torch.norm(noise_pred, dim=-1, keepdim=True)
noise_norm = torch.norm(comb_pred, dim=-1, keepdim=True)
noise_pred = comb_pred * (cond_norm / noise_norm)
# compute the previous noisy sample x_t -> x_t-1
latents_dtype = latents.dtype
latents = self.scheduler.step(noise_pred, t, latents, return_dict=False)[0]
if latents.dtype != latents_dtype:
if torch.backends.mps.is_available():
# some platforms (eg. apple mps) misbehave due to a pytorch bug: https://github.com/pytorch/pytorch/pull/99272
latents = latents.to(latents_dtype)
if callback_on_step_end is not None:
callback_kwargs = {}
for k in callback_on_step_end_tensor_inputs:
callback_kwargs[k] = locals()[k]
callback_outputs = callback_on_step_end(self, i, t, callback_kwargs)
latents = callback_outputs.pop("latents", latents)
prompt_embeds = callback_outputs.pop("prompt_embeds", prompt_embeds)
# call the callback, if provided
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0):
progress_bar.update()
if XLA_AVAILABLE:
xm.mark_step()
self._current_timestep = None
if output_type == "latent":
image = latents
else:
latents = self._unpack_latents(latents, height, width, self.vae_scale_factor)
latents = latents.to(self.vae.dtype)
latents_mean = (
torch.tensor(self.vae.config.latents_mean)
.view(1, self.vae.config.z_dim, 1, 1, 1)
.to(latents.device, latents.dtype)
)
latents_std = 1.0 / torch.tensor(self.vae.config.latents_std).view(1, self.vae.config.z_dim, 1, 1, 1).to(
latents.device, latents.dtype
)
latents = latents / latents_std + latents_mean
image = self.vae.decode(latents, return_dict=False)[0][:, :, 0]
image = self.image_processor.postprocess(image, output_type=output_type)
# Offload all models
self.maybe_free_model_hooks()
if not return_dict:
return (image,)
return QwenImagePipelineOutput(images=image)
@@ -296,9 +296,6 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if prompt_embeds is None:
prompt_embeds, prompt_embeds_mask = self._get_qwen_prompt_embeds(prompt, device)
prompt_embeds = prompt_embeds[:, :max_sequence_length]
prompt_embeds_mask = prompt_embeds_mask[:, :max_sequence_length]
_, seq_len, _ = prompt_embeds.shape
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
@@ -366,6 +363,21 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if max_sequence_length is not None and max_sequence_length > 1024:
raise ValueError(f"`max_sequence_length` cannot be greater than 1024 but is {max_sequence_length}")
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._prepare_latent_image_ids
def _prepare_latent_image_ids(batch_size, height, width, device, dtype):
latent_image_ids = torch.zeros(height, width, 3)
latent_image_ids[..., 1] = latent_image_ids[..., 1] + torch.arange(height)[:, None]
latent_image_ids[..., 2] = latent_image_ids[..., 2] + torch.arange(width)[None, :]
latent_image_id_height, latent_image_id_width, latent_image_id_channels = latent_image_ids.shape
latent_image_ids = latent_image_ids.reshape(
latent_image_id_height * latent_image_id_width, latent_image_id_channels
)
return latent_image_ids.to(device=device, dtype=dtype)
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._pack_latents
def _pack_latents(latents, batch_size, num_channels_latents, height, width):
@@ -453,7 +465,8 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
raise ValueError(f"Expected image dims 4 or 5, got {image.dim()}.")
if latents is not None:
return latents.to(device=device, dtype=dtype)
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents.to(device=device, dtype=dtype), latent_image_ids
image = image.to(device=device, dtype=dtype)
if image.shape[1] != self.latent_channels:
@@ -476,7 +489,9 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
latents = self.scheduler.scale_noise(image_latents, timestep, noise)
latents = self._pack_latents(latents, batch_size, num_channels_latents, height, width)
return latents
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents, latent_image_ids
@property
def guidance_scale(self):
@@ -698,7 +713,7 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
# 5. Prepare latent variables
num_channels_latents = self.transformer.config.in_channels // 4
latents = self.prepare_latents(
latents, latent_image_ids = self.prepare_latents(
init_image,
latent_timestep,
batch_size * num_images_per_prompt,
@@ -710,7 +725,7 @@ class QwenImageImg2ImgPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
generator,
latents,
)
img_shapes = [[(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)]] * batch_size
img_shapes = [(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)] * batch_size
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0)
self._num_timesteps = len(timesteps)
@@ -307,9 +307,6 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if prompt_embeds is None:
prompt_embeds, prompt_embeds_mask = self._get_qwen_prompt_embeds(prompt, device)
prompt_embeds = prompt_embeds[:, :max_sequence_length]
prompt_embeds_mask = prompt_embeds_mask[:, :max_sequence_length]
_, seq_len, _ = prompt_embeds.shape
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
@@ -393,6 +390,21 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
if max_sequence_length is not None and max_sequence_length > 1024:
raise ValueError(f"`max_sequence_length` cannot be greater than 1024 but is {max_sequence_length}")
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._prepare_latent_image_ids
def _prepare_latent_image_ids(batch_size, height, width, device, dtype):
latent_image_ids = torch.zeros(height, width, 3)
latent_image_ids[..., 1] = latent_image_ids[..., 1] + torch.arange(height)[:, None]
latent_image_ids[..., 2] = latent_image_ids[..., 2] + torch.arange(width)[None, :]
latent_image_id_height, latent_image_id_width, latent_image_id_channels = latent_image_ids.shape
latent_image_ids = latent_image_ids.reshape(
latent_image_id_height * latent_image_id_width, latent_image_id_channels
)
return latent_image_ids.to(device=device, dtype=dtype)
@staticmethod
# Copied from diffusers.pipelines.qwenimage.pipeline_qwenimage.QwenImagePipeline._pack_latents
def _pack_latents(latents, batch_size, num_channels_latents, height, width):
@@ -480,7 +492,8 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
raise ValueError(f"Expected image dims 4 or 5, got {image.dim()}.")
if latents is not None:
return latents.to(device=device, dtype=dtype)
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents.to(device=device, dtype=dtype), latent_image_ids
image = image.to(device=device, dtype=dtype)
if image.shape[1] != self.latent_channels:
@@ -511,7 +524,9 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
image_latents = self._pack_latents(image_latents, batch_size, num_channels_latents, height, width)
latents = self._pack_latents(latents, batch_size, num_channels_latents, height, width)
return latents, noise, image_latents
latent_image_ids = self._prepare_latent_image_ids(batch_size, height // 2, width // 2, device, dtype)
return latents, noise, image_latents, latent_image_ids
def prepare_mask_latents(
self,
@@ -844,7 +859,7 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
# 5. Prepare latent variables
num_channels_latents = self.transformer.config.in_channels // 4
latents, noise, image_latents = self.prepare_latents(
latents, noise, image_latents, latent_image_ids = self.prepare_latents(
init_image,
latent_timestep,
batch_size * num_images_per_prompt,
@@ -879,7 +894,7 @@ class QwenImageInpaintPipeline(DiffusionPipeline, QwenImageLoraLoaderMixin):
generator,
)
img_shapes = [[(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)]] * batch_size
img_shapes = [(1, height // self.vae_scale_factor // 2, width // self.vae_scale_factor // 2)] * batch_size
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0)
self._num_timesteps = len(timesteps)
@@ -525,7 +525,8 @@ class WanVACEPipeline(DiffusionPipeline, WanLoraLoaderMixin):
latents = retrieve_latents(self.vae.encode(video), generator, sample_mode="argmax").unbind(0)
latents = ((latents.float() - latents_mean) * latents_std).to(vae_dtype)
else:
mask = torch.where(mask > 0.5, 1.0, 0.0).to(dtype=vae_dtype)
mask = mask.to(dtype=vae_dtype)
mask = torch.where(mask > 0.5, 1.0, 0.0)
inactive = video * (1 - mask)
reactive = video * mask
inactive = retrieve_latents(self.vae.encode(inactive), generator, sample_mode="argmax")
-11
View File
@@ -209,17 +209,6 @@ class DiffusersQuantizer(ABC):
return model
def get_cuda_warm_up_factor(self):
"""
The factor to be used in `caching_allocator_warmup` to get the number of bytes to pre-allocate to warm up cuda.
A factor of 2 means we allocate all bytes in the empty model (since we allocate in fp16), a factor of 4 means
we allocate half the memory of the weights residing in the empty model, etc...
"""
# By default we return 4, i.e. half the model size (this corresponds to the case where the model is not
# really pre-processed, i.e. we do not have the info that weights are going to be 8 bits before actual
# weight loading)
return 4
def _dequantize(self, model):
raise NotImplementedError(
f"{self.quantization_config.quant_method} has no implementation of `dequantize`, please raise an issue on GitHub."
@@ -19,7 +19,6 @@ https://github.com/huggingface/transformers/blob/3a8eb74668e9c2cc563b2f5c62fac17
import importlib
import types
from fnmatch import fnmatch
from typing import TYPE_CHECKING, Any, Dict, List, Union
from packaging import version
@@ -279,31 +278,6 @@ class TorchAoHfQuantizer(DiffusersQuantizer):
module._parameters[tensor_name] = torch.nn.Parameter(param_value).to(device=target_device)
quantize_(module, self.quantization_config.get_apply_tensor_subclass())
def get_cuda_warm_up_factor(self):
"""
This factor is used in caching_allocator_warmup to determine how many bytes to pre-allocate for CUDA warmup.
- A factor of 2 means we pre-allocate the full memory footprint of the model.
- A factor of 4 means we pre-allocate half of that, and so on
However, when using TorchAO, calculating memory usage with param.numel() * param.element_size() doesn't give
the correct size for quantized weights (like int4 or int8) That's because TorchAO internally represents
quantized tensors using subtensors and metadata, and the reported element_size() still corresponds to the
torch_dtype not the actual bit-width of the quantized data.
To correct for this:
- Use a division factor of 8 for int4 weights
- Use a division factor of 4 for int8 weights
"""
# Original mapping for non-AOBaseConfig types
# For the uint types, this is a best guess. Once these types become more used
# we can look into their nuances.
map_to_target_dtype = {"int4_*": 8, "int8_*": 4, "uint*": 8, "float8*": 4}
quant_type = self.quantization_config.quant_type
for pattern, target_dtype in map_to_target_dtype.items():
if fnmatch(quant_type, pattern):
return target_dtype
raise ValueError(f"Unsupported quant_type: {quant_type!r}")
def _process_model_before_weight_loading(
self,
model: "ModelMixin",
+1 -1
View File
@@ -25,8 +25,8 @@ from .constants import (
DIFFUSERS_DYNAMIC_MODULE_NAME,
FLAX_WEIGHTS_NAME,
GGUF_FILE_EXTENSION,
HF_ENABLE_PARALLEL_LOADING,
HF_MODULES_CACHE,
HF_PARALLEL_LOADING_FLAG,
HUGGINGFACE_CO_RESOLVE_ENDPOINT,
MIN_PEFT_VERSION,
ONNX_EXTERNAL_WEIGHTS_NAME,
+1 -2
View File
@@ -44,8 +44,7 @@ DIFFUSERS_REQUEST_TIMEOUT = 60
DIFFUSERS_ATTN_BACKEND = os.getenv("DIFFUSERS_ATTN_BACKEND", "native")
DIFFUSERS_ATTN_CHECKS = os.getenv("DIFFUSERS_ATTN_CHECKS", "0") in ENV_VARS_TRUE_VALUES
DEFAULT_HF_PARALLEL_LOADING_WORKERS = 8
HF_ENABLE_PARALLEL_LOADING = os.environ.get("HF_ENABLE_PARALLEL_LOADING", "").upper() in ENV_VARS_TRUE_VALUES
DIFFUSERS_DISABLE_REMOTE_CODE = os.getenv("DIFFUSERS_DISABLE_REMOTE_CODE", "false").lower() in ENV_VARS_TRUE_VALUES
HF_PARALLEL_LOADING_FLAG = "HF_ENABLE_PARALLEL_LOADING"
# Below should be `True` if the current version of `peft` and `transformers` are compatible with
# PEFT backend. Will automatically fall back to PEFT backend if the correct versions of the libraries are
@@ -1742,21 +1742,6 @@ class PixArtSigmaPipeline(metaclass=DummyObject):
requires_backends(cls, ["torch", "transformers"])
class QwenImageEditPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
def __init__(self, *args, **kwargs):
requires_backends(self, ["torch", "transformers"])
@classmethod
def from_config(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
@classmethod
def from_pretrained(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
class QwenImageImg2ImgPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
+44 -17
View File
@@ -20,6 +20,7 @@ import json
import os
import re
import shutil
import signal
import sys
import threading
from pathlib import Path
@@ -33,7 +34,6 @@ from packaging import version
from .. import __version__
from . import DIFFUSERS_DYNAMIC_MODULE_NAME, HF_MODULES_CACHE, logging
from .constants import DIFFUSERS_DISABLE_REMOTE_CODE
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
@@ -159,25 +159,52 @@ def check_imports(filename):
return get_relative_imports(filename)
def _raise_timeout_error(signum, frame):
raise ValueError(
"Loading this model requires you to execute custom code contained in the model repository on your local "
"machine. Please set the option `trust_remote_code=True` to permit loading of this model."
)
def resolve_trust_remote_code(trust_remote_code, model_name, has_remote_code):
trust_remote_code = trust_remote_code and not DIFFUSERS_DISABLE_REMOTE_CODE
if DIFFUSERS_DISABLE_REMOTE_CODE:
logger.warning(
"Downloading remote code is disabled globally via the DIFFUSERS_DISABLE_REMOTE_CODE environment variable. Ignoring `trust_remote_code`."
)
if trust_remote_code is None:
if has_remote_code and TIME_OUT_REMOTE_CODE > 0:
prev_sig_handler = None
try:
prev_sig_handler = signal.signal(signal.SIGALRM, _raise_timeout_error)
signal.alarm(TIME_OUT_REMOTE_CODE)
while trust_remote_code is None:
answer = input(
f"The repository for {model_name} contains custom code which must be executed to correctly "
f"load the model. You can inspect the repository content at https://hf.co/{model_name}.\n"
f"You can avoid this prompt in future by passing the argument `trust_remote_code=True`.\n\n"
f"Do you wish to run the custom code? [y/N] "
)
if answer.lower() in ["yes", "y", "1"]:
trust_remote_code = True
elif answer.lower() in ["no", "n", "0", ""]:
trust_remote_code = False
signal.alarm(0)
except Exception:
# OS which does not support signal.SIGALRM
raise ValueError(
f"The repository for {model_name} contains custom code which must be executed to correctly "
f"load the model. You can inspect the repository content at https://hf.co/{model_name}.\n"
f"Please pass the argument `trust_remote_code=True` to allow custom code to be run."
)
finally:
if prev_sig_handler is not None:
signal.signal(signal.SIGALRM, prev_sig_handler)
signal.alarm(0)
elif has_remote_code:
# For the CI which puts the timeout at 0
_raise_timeout_error(None, None)
if has_remote_code and not trust_remote_code:
error_msg = f"The repository for {model_name} contains custom code. "
error_msg += (
"Downloading remote code is disabled globally via the DIFFUSERS_DISABLE_REMOTE_CODE environment variable."
if DIFFUSERS_DISABLE_REMOTE_CODE
else "Pass `trust_remote_code=True` to allow loading remote code modules."
)
raise ValueError(error_msg)
elif has_remote_code and trust_remote_code:
logger.warning(
f"`trust_remote_code` is enabled. Downloading code from {model_name}. Please ensure you trust the contents of this repository"
raise ValueError(
f"Loading {model_name} requires you to execute the configuration file in that"
" repo on your local machine. Make sure you have read the code there to avoid malicious use, then"
" set the option `trust_remote_code=True` to remove this error."
)
return trust_remote_code
-2
View File
@@ -15,7 +15,6 @@
PyTorch utilities: Utilities related to PyTorch
"""
import functools
from typing import List, Optional, Tuple, Union
from . import logging
@@ -169,7 +168,6 @@ def get_torch_cuda_device_capability():
return None
@functools.lru_cache
def get_device():
if torch.cuda.is_available():
return "cuda"
+6
View File
@@ -18,6 +18,7 @@ import tempfile
import unittest
import numpy as np
import pytest
import safetensors.torch
import torch
from PIL import Image
@@ -159,6 +160,11 @@ class WanVACELoRATests(unittest.TestCase, PeftLoraLoaderMixinTests):
def test_simple_inference_with_text_lora_save_load(self):
pass
@pytest.mark.xfail(
condition=True,
reason="RuntimeError: Input type (float) and bias type (c10::BFloat16) should be the same",
strict=True,
)
def test_layerwise_casting_inference_denoiser(self):
super().test_layerwise_casting_inference_denoiser()
@@ -15,7 +15,6 @@
import unittest
import pytest
import torch
from diffusers import QwenImageTransformer2DModel
@@ -100,7 +99,3 @@ class QwenImageTransformerCompileTests(TorchCompileTesterMixin, unittest.TestCas
def prepare_dummy_input(self, height, width):
return QwenImageTransformerTests().prepare_dummy_input(height=height, width=width)
@pytest.mark.xfail(condition=True, reason="RoPE needs to be revisited.", strict=True)
def test_torch_compile_recompilation_and_graph_break(self):
super().test_torch_compile_recompilation_and_graph_break()
@@ -1,243 +0,0 @@
# Copyright 2025 The HuggingFace Team.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import unittest
import numpy as np
import pytest
import torch
from PIL import Image
from transformers import Qwen2_5_VLConfig, Qwen2_5_VLForConditionalGeneration, Qwen2Tokenizer, Qwen2VLProcessor
from diffusers import (
AutoencoderKLQwenImage,
FlowMatchEulerDiscreteScheduler,
QwenImageEditPipeline,
QwenImageTransformer2DModel,
)
from diffusers.utils.testing_utils import enable_full_determinism, torch_device
from ..pipeline_params import TEXT_TO_IMAGE_PARAMS
from ..test_pipelines_common import PipelineTesterMixin, to_np
enable_full_determinism()
class QwenImageEditPipelineFastTests(PipelineTesterMixin, unittest.TestCase):
pipeline_class = QwenImageEditPipeline
params = TEXT_TO_IMAGE_PARAMS - {"cross_attention_kwargs"}
batch_params = frozenset(["prompt", "image"])
image_params = frozenset(["image"])
image_latents_params = frozenset(["latents"])
required_optional_params = frozenset(
[
"num_inference_steps",
"generator",
"latents",
"return_dict",
"callback_on_step_end",
"callback_on_step_end_tensor_inputs",
]
)
supports_dduf = False
test_xformers_attention = False
test_layerwise_casting = True
test_group_offloading = True
def get_dummy_components(self):
tiny_ckpt_id = "hf-internal-testing/tiny-random-Qwen2VLForConditionalGeneration"
torch.manual_seed(0)
transformer = QwenImageTransformer2DModel(
patch_size=2,
in_channels=16,
out_channels=4,
num_layers=2,
attention_head_dim=16,
num_attention_heads=3,
joint_attention_dim=16,
guidance_embeds=False,
axes_dims_rope=(8, 4, 4),
)
torch.manual_seed(0)
z_dim = 4
vae = AutoencoderKLQwenImage(
base_dim=z_dim * 6,
z_dim=z_dim,
dim_mult=[1, 2, 4],
num_res_blocks=1,
temperal_downsample=[False, True],
latents_mean=[0.0] * z_dim,
latents_std=[1.0] * z_dim,
)
torch.manual_seed(0)
scheduler = FlowMatchEulerDiscreteScheduler()
torch.manual_seed(0)
config = Qwen2_5_VLConfig(
text_config={
"hidden_size": 16,
"intermediate_size": 16,
"num_hidden_layers": 2,
"num_attention_heads": 2,
"num_key_value_heads": 2,
"rope_scaling": {
"mrope_section": [1, 1, 2],
"rope_type": "default",
"type": "default",
},
"rope_theta": 1000000.0,
},
vision_config={
"depth": 2,
"hidden_size": 16,
"intermediate_size": 16,
"num_heads": 2,
"out_hidden_size": 16,
},
hidden_size=16,
vocab_size=152064,
vision_end_token_id=151653,
vision_start_token_id=151652,
vision_token_id=151654,
)
text_encoder = Qwen2_5_VLForConditionalGeneration(config)
tokenizer = Qwen2Tokenizer.from_pretrained(tiny_ckpt_id)
components = {
"transformer": transformer,
"vae": vae,
"scheduler": scheduler,
"text_encoder": text_encoder,
"tokenizer": tokenizer,
"processor": Qwen2VLProcessor.from_pretrained(tiny_ckpt_id),
}
return components
def get_dummy_inputs(self, device, seed=0):
if str(device).startswith("mps"):
generator = torch.manual_seed(seed)
else:
generator = torch.Generator(device=device).manual_seed(seed)
inputs = {
"prompt": "dance monkey",
"image": Image.new("RGB", (32, 32)),
"negative_prompt": "bad quality",
"generator": generator,
"num_inference_steps": 2,
"true_cfg_scale": 1.0,
"height": 32,
"width": 32,
"max_sequence_length": 16,
"output_type": "pt",
}
return inputs
def test_inference(self):
device = "cpu"
components = self.get_dummy_components()
pipe = self.pipeline_class(**components)
pipe.to(device)
pipe.set_progress_bar_config(disable=None)
inputs = self.get_dummy_inputs(device)
image = pipe(**inputs).images
generated_image = image[0]
self.assertEqual(generated_image.shape, (3, 32, 32))
# fmt: off
expected_slice = torch.tensor([[0.5637, 0.6341, 0.6001, 0.5620, 0.5794, 0.5498, 0.5757, 0.6389, 0.4174, 0.3597, 0.5649, 0.4894, 0.4969, 0.5255, 0.4083, 0.4986]])
# fmt: on
generated_slice = generated_image.flatten()
generated_slice = torch.cat([generated_slice[:8], generated_slice[-8:]])
self.assertTrue(torch.allclose(generated_slice, expected_slice, atol=1e-3))
def test_inference_batch_single_identical(self):
self._test_inference_batch_single_identical(batch_size=3, expected_max_diff=1e-1)
def test_attention_slicing_forward_pass(
self, test_max_difference=True, test_mean_pixel_difference=True, expected_max_diff=1e-3
):
if not self.test_attention_slicing:
return
components = self.get_dummy_components()
pipe = self.pipeline_class(**components)
for component in pipe.components.values():
if hasattr(component, "set_default_attn_processor"):
component.set_default_attn_processor()
pipe.to(torch_device)
pipe.set_progress_bar_config(disable=None)
generator_device = "cpu"
inputs = self.get_dummy_inputs(generator_device)
output_without_slicing = pipe(**inputs)[0]
pipe.enable_attention_slicing(slice_size=1)
inputs = self.get_dummy_inputs(generator_device)
output_with_slicing1 = pipe(**inputs)[0]
pipe.enable_attention_slicing(slice_size=2)
inputs = self.get_dummy_inputs(generator_device)
output_with_slicing2 = pipe(**inputs)[0]
if test_max_difference:
max_diff1 = np.abs(to_np(output_with_slicing1) - to_np(output_without_slicing)).max()
max_diff2 = np.abs(to_np(output_with_slicing2) - to_np(output_without_slicing)).max()
self.assertLess(
max(max_diff1, max_diff2),
expected_max_diff,
"Attention slicing should not affect the inference results",
)
def test_vae_tiling(self, expected_diff_max: float = 0.2):
generator_device = "cpu"
components = self.get_dummy_components()
pipe = self.pipeline_class(**components)
pipe.to("cpu")
pipe.set_progress_bar_config(disable=None)
# Without tiling
inputs = self.get_dummy_inputs(generator_device)
inputs["height"] = inputs["width"] = 128
output_without_tiling = pipe(**inputs)[0]
# With tiling
pipe.vae.enable_tiling(
tile_sample_min_height=96,
tile_sample_min_width=96,
tile_sample_stride_height=64,
tile_sample_stride_width=64,
)
inputs = self.get_dummy_inputs(generator_device)
inputs["height"] = inputs["width"] = 128
output_with_tiling = pipe(**inputs)[0]
self.assertLess(
(to_np(output_without_tiling) - to_np(output_with_tiling)).max(),
expected_diff_max,
"VAE tiling should not affect the inference results",
)
@pytest.mark.xfail(condition=True, reason="Preconfigured embeddings need to be revisited.", strict=True)
def test_encode_prompt_works_in_isolation(self, extra_required_param_value_dict=None, atol=1e-4, rtol=1e-4):
super().test_encode_prompt_works_in_isolation(extra_required_param_value_dict, atol, rtol)
-23
View File
@@ -2339,29 +2339,6 @@ class PipelineTesterMixin:
f"Component '{name}' has dtype {component.dtype} but expected {expected_dtype}",
)
@require_torch_accelerator
def test_pipeline_with_accelerator_device_map(self, expected_max_difference=1e-4):
components = self.get_dummy_components()
pipe = self.pipeline_class(**components)
pipe = pipe.to(torch_device)
pipe.set_progress_bar_config(disable=None)
torch.manual_seed(0)
inputs = self.get_dummy_inputs(torch_device)
inputs["generator"] = torch.manual_seed(0)
out = pipe(**inputs)[0]
with tempfile.TemporaryDirectory() as tmpdir:
pipe.save_pretrained(tmpdir)
loaded_pipe = self.pipeline_class.from_pretrained(tmpdir, device_map=torch_device)
for component in loaded_pipe.components.values():
if hasattr(component, "set_default_attn_processor"):
component.set_default_attn_processor()
inputs["generator"] = torch.manual_seed(0)
loaded_out = loaded_pipe(**inputs)[0]
max_diff = np.abs(to_np(out) - to_np(loaded_out)).max()
self.assertLess(max_diff, expected_max_difference)
@is_staging_test
class PipelinePushToHubTester(unittest.TestCase):
+27 -27
View File
@@ -304,7 +304,7 @@ class FluxGGUFSingleFileTests(GGUFSingleFileTesterMixin, unittest.TestCase):
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
config="black-forest-labs/FLUX.1-dev",
)
model.to(torch_device)
model.to("cuda")
model(**self.get_dummy_inputs())
@@ -360,33 +360,33 @@ class SD35LargeGGUFSingleFileTests(GGUFSingleFileTesterMixin, unittest.TestCase)
{
("xpu", 3): np.array(
[
0.1953125,
0.3125,
0.31445312,
0.13085938,
0.30664062,
0.29296875,
0.11523438,
0.2890625,
0.28320312,
0.16601562,
0.3046875,
0.328125,
0.140625,
0.31640625,
0.32421875,
0.12304688,
0.3046875,
0.3046875,
0.17578125,
0.3359375,
0.3203125,
0.16601562,
0.34375,
0.31640625,
0.15429688,
0.328125,
0.16210938,
0.2734375,
0.27734375,
0.109375,
0.27148438,
0.2578125,
0.1015625,
0.2578125,
0.2578125,
0.14453125,
0.26953125,
0.29492188,
0.12890625,
0.28710938,
0.30078125,
0.11132812,
0.27734375,
0.27929688,
0.15625,
0.31054688,
0.296875,
0.15234375,
0.3203125,
0.29492188,
0.140625,
0.3046875,
0.28515625,
]
),
("cuda", 7): np.array(