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Author SHA1 Message Date
Patrick von Platen 4719b8f5f9 Release: v0.23.1 2023-11-16 15:43:52 +01:00
Patrick von Platen a53c17b698 correct installs 2023-11-16 15:41:23 +01:00
Patrick von Platen 41d6ca0058 correct peft 2023-11-16 15:23:45 +01:00
Patrick von Platen a28a77bfe9 put peft in requirements 2023-11-16 15:20:08 +01:00
sayakpaul fbb8b34716 Release: v0.23.0 2023-11-09 21:55:21 +05:30
Sayak Paul bc2ba004c6 [LCM] add: locm docs. (#5723)
* add: locm docs.

* correct path

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* up

* add

---------

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-09 17:21:26 +01:00
Patrick von Platen 3d7eaf83d7 LCM Add Tests (#5707)
* lcm add tests

* uP

* Fix all

* uP

* Add

* all

* uP

* uP

* uP

* uP

* uP

* uP

* uP
2023-11-09 15:45:11 +01:00
Patrick von Platen bf406ea886 Correct consist dec (#5722)
* uP

* Update src/diffusers/models/consistency_decoder_vae.py

* uP

* uP
2023-11-09 13:10:24 +01:00
Will Berman 2fd46405cd consistency decoder (#5694)
* consistency decoder

* rename

* Apply suggestions from code review

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update src/diffusers/pipelines/consistency_models/pipeline_consistency_models.py

* uP

* Apply suggestions from code review

* uP

* uP

* uP

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-11-09 12:21:41 +01:00
Dhruv Nair 43346adc1f Install accelerate from PyPI in PR test runner (#5721)
install acclerate from pypi
2023-11-09 15:38:00 +05:30
takuoko 6110d7c95f [Bugfix] fix error of peft lora when xformers enabled (#5697)
* bugfix peft lor

* Apply suggestions from code review

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-08 22:55:36 +01:00
Dhruv Nair 65ef7a0c5c Fix prompt bug in AnimateDiff (#5702)
* fix prompt bug

* add test
2023-11-08 21:49:09 +01:00
apolinário 6e68c71503 Add adapter fusing + PEFT to the docs (#5662)
* Add adapter fusing + PEFT to the docs

* Update docs/source/en/tutorials/using_peft_for_inference.md

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update docs/source/en/tutorials/using_peft_for_inference.md

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update docs/source/en/tutorials/using_peft_for_inference.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/tutorials/using_peft_for_inference.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/tutorials/using_peft_for_inference.md

* Update docs/source/en/tutorials/using_peft_for_inference.md

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2023-11-08 18:26:53 +01:00
Patrick von Platen 17528afcba Fix styling issues (#5699)
* up

* up

* up

* Empty-Commit

* fix keyword argument call.

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-11-08 17:49:06 +05:30
Sayak Paul 78be400761 [PixArt-Alpha] fix mask feature condition. (#5695)
* fix mask feature condition.

* debug

* remove identical test

* set correct

* Empty-Commit
2023-11-08 17:42:46 +05:30
Patrick von Platen c803a8f8c0 [LCM] Fix img2img (#5698)
* [LCM] Fix img2img

* make fix-copies

* make fix-copies

* make fix-copies

* up
2023-11-08 11:51:46 +01:00
Philipp Hasper d384265df7 Fixed is_safetensors_compatible() handling of windows path separators (#5650)
Closes #4665
2023-11-08 11:51:15 +01:00
Kirill 11c125667b Fix the misaligned pipeline usage in dreamshaper docstrings (#5700)
Fix the misaligned pipeline usage
2023-11-08 11:49:03 +01:00
YiYi Xu 69996938cf speed up Shap-E fast test (#5686)
skip rendering

Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-11-08 11:43:20 +01:00
Chi 9ae90593c0 Replacing the nn.Mish activation function with a get_activation function. (#5651)
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.

* Update src/diffusers/models/unet_2d_blocks.py

This changes suggest by maintener.

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update src/diffusers/models/unet_2d_blocks.py

Add suggested text

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update unet_2d_blocks.py

I changed the Parameter to Args text.

* Update unet_2d_blocks.py

proper indentation set in this file.

* Update unet_2d_blocks.py

a little bit of change in the act_fun argument line.

* I run the black command to reformat style in the code

* Update unet_2d_blocks.py

similar doc-string add to have in the original diffusion repository.

* I removed the dummy variable defined in both the encoder and decoder.

* Now, I run black package to reformat my file

* Remove the redundant line from the adapter.py file.

* Black package using to reformated my file

* Replacing the nn.Mish activation function with a get_activation function allows developers to more easily choose the right activation function for their task. Additionally, removing redundant variables can improve code readability and maintainability.

* I try to fix this: Fast tests for PRs / Fast PyTorch Models & Schedulers CPU tests (pull_request)

* Update src/diffusers/models/resnet.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
2023-11-08 07:58:21 +05:30
M. Tolga Cangöz 7942bb8dc2 [Docs] Fix typos, improve, update at Using Diffusers' Task page (#5611)
* Fix typos, improve, update; kandinsky doesn't want fp16 due to deprecation; ogkalu and kohbanye don't have safetensor; add make_image_grid for better visualization

* Update inpaint.md

* Remove erronous Space

* Update docs/source/en/using-diffusers/conditional_image_generation.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update img2img.md

* load_image() already converts to RGB

* Update depth2img.md

* Update img2img.md

* Update inpaint.md

---------

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2023-11-07 15:42:20 -08:00
dg845 aab6de22c3 Improve LCMScheduler (#5681)
* Refactor LCMScheduler.step such that prev_sample == denoised at the last timestep in the schedule.

* Make timestep scaling when calculating boundary conditions configurable.

* Reparameterize timestep_scaling to be a multiplicative rather than division scaling.

* make style

* fix dtype conversion

* make style

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-07 18:48:18 +01:00
Sayak Paul 1dc231d14a [PixArt-Alpha] Support non-square images (#5672)
* debug

* support non-square images

* add: test

* fix: test

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-07 18:21:33 +01:00
Sayak Paul 84cd9e8d01 Make sure DDPM and diffusers can be used without Transformers (#5668)
* fix: import bug

* fix

* fix

* fix import utils for lcm

* fix: pixart alpha init

* Fix

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-07 17:38:12 +01:00
Sayak Paul a8523bffa8 [PixArt-Alpha] fix mask_feature so that precomputed embeddings work with a batch size > 1 (#5677)
* fix embeds

* remove todo

* add: test

* better name
2023-11-07 17:12:47 +01:00
Dhruv Nair 97c8199dbb Explicit torch/flax dependency check (#5673)
* explicit torch dependency check

* update

* update

* update
2023-11-07 16:38:20 +01:00
Dhruv Nair 414d7c4991 Fix Basic Transformer Block (#5683)
* fix

* Update src/diffusers/models/attention.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-07 16:36:49 +01:00
Dhruv Nair 8ca179a0a9 Update free model hooks (#5680)
update free model hooks
2023-11-07 20:50:57 +05:30
Dhruv Nair 71f56c771a Model tests xformers fixes (#5679)
* fix model xformers test

* update
2023-11-07 20:50:41 +05:30
Dhruv Nair 6a89a6c93a Update custom diffusion attn processor (#5663)
update custom diffusion attn processor
2023-11-07 12:46:38 +05:30
117 changed files with 4245 additions and 970 deletions
+1 -2
View File
@@ -28,6 +28,5 @@ jobs:
pip install pytest
- name: Check for soft dependencies
run: |
python utils/print_env.py
python -c "from diffusers import UNet2DConditionModel; print(type(UNet2DConditionModel))"
pytest tests/others/test_dependencies.py
@@ -0,0 +1,34 @@
name: Run Flax dependency tests
on:
pull_request:
branches:
- main
push:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
check_flax_dependencies:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -e .
pip install "jax[cpu]>=0.2.16,!=0.3.2"
pip install "flax>=0.4.1"
pip install "jaxlib>=0.1.65"
pip install pytest
- name: Check for soft dependencies
run: |
pytest tests/others/test_dependencies.py
+2 -2
View File
@@ -72,7 +72,7 @@ jobs:
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
python -m pip install accelerate
- name: Environment
run: |
@@ -115,7 +115,7 @@ jobs:
run: |
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples/test_examples.py
examples/test_examples.py
- name: Failure short reports
if: ${{ failure() }}
@@ -0,0 +1,32 @@
name: Run Torch dependency tests
on:
pull_request:
branches:
- main
push:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
check_torch_dependencies:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -e .
pip install torch torchvision torchaudio
pip install pytest
- name: Check for soft dependencies
run: |
pytest tests/others/test_dependencies.py
+6
View File
@@ -72,6 +72,8 @@
title: Overview
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/lcm
title: Latent Consistency Models
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/controlnet
@@ -200,6 +202,8 @@
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
title: Transformer2D
- local: api/models/transformer_temporal
@@ -344,6 +348,8 @@
title: Overview
- local: api/schedulers/cm_stochastic_iterative
title: CMStochasticIterativeScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/ddim_inverse
title: DDIMInverseScheduler
- local: api/schedulers/ddim
@@ -0,0 +1,18 @@
# Consistency Decoder
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistencydecoder](https://github.com/openai/consistencydecoder).
<Tip warning={true}>
Inference is only supported for 2 iterations as of now.
</Tip>
The pipeline could not have been contributed without the help of [madebyollin](https://github.com/madebyollin) and [mrsteyk](https://github.com/mrsteyk) from [this issue](https://github.com/openai/consistencydecoder/issues/1).
## ConsistencyDecoderVAE
[[autodoc]] ConsistencyDecoderVAE
- all
- decode
@@ -0,0 +1,9 @@
# ConsistencyDecoderScheduler
This scheduler is a part of the [`ConsistencyDecoderPipeline`] and was introduced in [DALL-E 3](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models).
## ConsistencyDecoderScheduler
[[autodoc]] schedulers.scheduling_consistency_decoder.ConsistencyDecoderScheduler
@@ -12,9 +12,9 @@ specific language governing permissions and limitations under the License.
[[open-in-colab]]
# Inference with PEFT
# Load LoRAs for inference
There are many adapters trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images. With the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers, it is really easy to load and manage adapters for inference. In this guide, you'll learn how to use different adapters with [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) for inference.
There are many adapters (with LoRAs being the most common type) trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images. With the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers, it is really easy to load and manage adapters for inference. In this guide, you'll learn how to use different adapters with [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) for inference.
Throughout this guide, you'll use LoRA as the main adapter technique, so we'll use the terms LoRA and adapter interchangeably. You should have some familiarity with LoRA, and if you don't, we welcome you to check out the [LoRA guide](https://huggingface.co/docs/peft/conceptual_guides/lora).
@@ -22,9 +22,8 @@ Let's first install all the required libraries.
```bash
!pip install -q transformers accelerate
# Will be updated once the stable releases are done.
!pip install -q git+https://github.com/huggingface/peft.git
!pip install -q git+https://github.com/huggingface/diffusers.git
!pip install peft
!pip install diffusers
```
Now, let's load a pipeline with a SDXL checkpoint:
@@ -165,3 +164,22 @@ list_adapters_component_wise = pipe.get_list_adapters()
list_adapters_component_wise
{"text_encoder": ["toy", "pixel"], "unet": ["toy", "pixel"], "text_encoder_2": ["toy", "pixel"]}
```
## Fusing adapters into the model
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~diffusers.loaders.LoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora()
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
# Gets the Unet back to the original state
pipe.unfuse_lora()
```
@@ -30,6 +30,7 @@ You can generate images from a prompt in 🤗 Diffusers in two steps:
```py
from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16"
@@ -42,6 +43,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
image = pipeline(
"stained glass of darth vader, backlight, centered composition, masterpiece, photorealistic, 8k"
).images[0]
image
```
<div class="flex justify-center">
@@ -65,6 +67,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
).to("cuda")
generator = torch.Generator("cuda").manual_seed(31)
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", generator=generator).images[0]
image
```
### Stable Diffusion XL
@@ -80,6 +83,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
).to("cuda")
generator = torch.Generator("cuda").manual_seed(31)
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", generator=generator).images[0]
image
```
### Kandinsky 2.2
@@ -93,15 +97,16 @@ from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, variant="fp16"
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16
).to("cuda")
generator = torch.Generator("cuda").manual_seed(31)
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", generator=generator).images[0]
image
```
### ControlNet
ControlNet are auxiliary models or adapters that are finetuned on top of text-to-image models, such as [Stable Diffusion V1.5](https://huggingface.co/runwayml/stable-diffusion-v1-5). Using ControlNet models in combination with text-to-image models offers diverse options for more explicit control over how to generate an image. With ControlNet's, you add an additional conditioning input image to the model. For example, if you provide an image of a human pose (usually represented as multiple keypoints that are connected into a skeleton) as a conditioning input, the model generates an image that follows the pose of the image. Check out the more in-depth [ControlNet](controlnet) guide to learn more about other conditioning inputs and how to use them.
ControlNet models are auxiliary models or adapters that are finetuned on top of text-to-image models, such as [Stable Diffusion v1.5](https://huggingface.co/runwayml/stable-diffusion-v1-5). Using ControlNet models in combination with text-to-image models offers diverse options for more explicit control over how to generate an image. With ControlNet, you add an additional conditioning input image to the model. For example, if you provide an image of a human pose (usually represented as multiple keypoints that are connected into a skeleton) as a conditioning input, the model generates an image that follows the pose of the image. Check out the more in-depth [ControlNet](controlnet) guide to learn more about other conditioning inputs and how to use them.
In this example, let's condition the ControlNet with a human pose estimation image. Load the ControlNet model pretrained on human pose estimations:
@@ -124,6 +129,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
).to("cuda")
generator = torch.Generator("cuda").manual_seed(31)
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", image=pose_image, generator=generator).images[0]
image
```
<div class="flex flex-row gap-4">
@@ -163,6 +169,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
image = pipeline(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", height=768, width=512
).images[0]
image
```
<div class="flex justify-center">
@@ -171,7 +178,7 @@ image = pipeline(
<Tip warning={true}>
Other models may have different default image sizes depending on the image size's in the training dataset. For example, SDXL's default image size is 1024x1024 and using lower `height` and `width` values may result in lower quality images. Make sure you check the model's API reference first!
Other models may have different default image sizes depending on the image sizes in the training dataset. For example, SDXL's default image size is 1024x1024 and using lower `height` and `width` values may result in lower quality images. Make sure you check the model's API reference first!
</Tip>
@@ -189,6 +196,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
image = pipeline(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", guidance_scale=3.5
).images[0]
image
```
<div class="flex flex-row gap-4">
@@ -221,16 +229,17 @@ image = pipeline(
prompt="Astronaut in a jungle, cold color palette, muted colors, detailed, 8k",
negative_prompt="ugly, deformed, disfigured, poor details, bad anatomy",
).images[0]
image
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/text2img-neg-prompt-1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative prompt = "ugly, deformed, disfigured, poor details, bad anatomy"</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative_prompt = "ugly, deformed, disfigured, poor details, bad anatomy"</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/text2img-neg-prompt-2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative prompt = "astronaut"</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative_prompt = "astronaut"</figcaption>
</div>
</div>
@@ -252,6 +261,7 @@ image = pipeline(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k",
generator=generator,
).images[0]
image
```
## Control image generation
@@ -278,14 +288,14 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16
).to("cuda")
image = pipeline(
prompt_emebds=prompt_embeds, # generated from Compel
prompt_embeds=prompt_embeds, # generated from Compel
negative_prompt_embeds=negative_prompt_embeds, # generated from Compel
).images[0]
```
### ControlNet
As you saw in the [ControlNet](#controlnet) section, these models offer a more flexible and accurate way to generate images by incorporating an additional conditioning image input. Each ControlNet model is pretrained on a particular type of conditioning image to generate new images that resemble it. For example, if you take a ControlNet pretrained on depth maps, you can give the model a depth map as a conditioning input and it'll generate an image that preserves the spatial information in it. This is quicker and easier than specifying the depth information in a prompt. You can even combine multiple conditioning inputs with a [MultiControlNet](controlnet#multicontrolnet)!
As you saw in the [ControlNet](#controlnet) section, these models offer a more flexible and accurate way to generate images by incorporating an additional conditioning image input. Each ControlNet model is pretrained on a particular type of conditioning image to generate new images that resemble it. For example, if you take a ControlNet model pretrained on depth maps, you can give the model a depth map as a conditioning input and it'll generate an image that preserves the spatial information in it. This is quicker and easier than specifying the depth information in a prompt. You can even combine multiple conditioning inputs with a [MultiControlNet](controlnet#multicontrolnet)!
There are many types of conditioning inputs you can use, and 🤗 Diffusers supports ControlNet for Stable Diffusion and SDXL models. Take a look at the more comprehensive [ControlNet](controlnet) guide to learn how you can use these models.
@@ -300,7 +310,7 @@ from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16").to("cuda")
pipeline.unet = torch.compile(pipeline.unet, mode="reduce-overheard", fullgraph=True)
pipeline.unet = torch.compile(pipeline.unet, mode="reduce-overhead", fullgraph=True)
```
For more tips on how to optimize your code to save memory and speed up inference, read the [Memory and speed](../optimization/fp16) and [Torch 2.0](../optimization/torch2.0) guides.
For more tips on how to optimize your code to save memory and speed up inference, read the [Memory and speed](../optimization/fp16) and [Torch 2.0](../optimization/torch2.0) guides.
+6 -17
View File
@@ -20,12 +20,10 @@ Start by creating an instance of the [`StableDiffusionDepth2ImgPipeline`]:
```python
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline
from diffusers.utils import load_image, make_image_grid
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
pipeline = StableDiffusionDepth2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-depth",
torch_dtype=torch.float16,
use_safetensors=True,
@@ -36,22 +34,13 @@ Now pass your prompt to the pipeline. You can also pass a `negative_prompt` to p
```python
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)
init_image = load_image(url)
prompt = "two tigers"
n_prompt = "bad, deformed, ugly, bad anatomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_prompt, strength=0.7).images[0]
image
negative_prompt = "bad, deformed, ugly, bad anatomy"
image = pipeline(prompt=prompt, image=init_image, negative_prompt=negative_prompt, strength=0.7).images[0]
make_image_grid([init_image, image], rows=1, cols=2)
```
| Input | Output |
|---------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------|
| <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/coco-cats.png" width="500"/> | <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/depth2img-tigers.png" width="500"/> |
Play around with the Spaces below and see if you notice a difference between generated images with and without a depth map!
<iframe
src="https://radames-stable-diffusion-depth2img.hf.space"
frameborder="0"
width="850"
height="500"
></iframe>
+78 -79
View File
@@ -21,13 +21,15 @@ With 🤗 Diffusers, this is as easy as 1-2-3:
1. Load a checkpoint into the [`AutoPipelineForImage2Image`] class; this pipeline automatically handles loading the correct pipeline class based on the checkpoint:
```py
import torch
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForImage2Image.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
```
@@ -48,7 +50,7 @@ init_image = load_image("https://huggingface.co/datasets/huggingface/documentati
```py
prompt = "cat wizard, gandalf, lord of the rings, detailed, fantasy, cute, adorable, Pixar, Disney, 8k"
image = pipeline(prompt, image=init_image).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex gap-4">
@@ -72,27 +74,25 @@ Stable Diffusion v1.5 is a latent diffusion model initialized from an earlier ch
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
# pass prompt and image to pipeline
image = pipeline(prompt, image=init_image).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex gap-4">
@@ -112,27 +112,25 @@ SDXL is a more powerful version of the Stable Diffusion model. It uses a larger
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-sdxl-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
# pass prompt and image to pipeline
image = pipeline(prompt, image=init_image, strength=0.5).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex gap-4">
@@ -154,27 +152,25 @@ The simplest way to use Kandinsky 2.2 is:
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
# pass prompt and image to pipeline
image = pipeline(prompt, image=init_image).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex gap-4">
@@ -199,32 +195,29 @@ There are several important parameters you can configure in the pipeline that'll
- 📈 a higher `strength` value gives the model more "creativity" to generate an image that's different from the initial image; a `strength` value of 1.0 means the initial image is more or less ignored
- 📉 a lower `strength` value means the generated image is more similar to the initial image
The `strength` and `num_inference_steps` parameter are related because `strength` determines the number of noise steps to add. For example, if the `num_inference_steps` is 50 and `strength` is 0.8, then this means adding 40 (50 * 0.8) steps of noise to the initial image and then denoising for 40 steps to get the newly generated image.
The `strength` and `num_inference_steps` parameters are related because `strength` determines the number of noise steps to add. For example, if the `num_inference_steps` is 50 and `strength` is 0.8, then this means adding 40 (50 * 0.8) steps of noise to the initial image and then denoising for 40 steps to get the newly generated image.
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = init_image
# pass prompt and image to pipeline
image = pipeline(prompt, image=init_image, strength=0.8).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex flex-row gap-4">
@@ -250,27 +243,25 @@ You can combine `guidance_scale` with `strength` for even more precise control o
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
# pass prompt and image to pipeline
image = pipeline(prompt, image=init_image, guidance_scale=8.0).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex flex-row gap-4">
@@ -294,38 +285,36 @@ A negative prompt conditions the model to *not* include things in an image, and
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
negative_prompt = "ugly, deformed, disfigured, poor details, bad anatomy"
# pass prompt and image to pipeline
image = pipeline(prompt, negative_prompt=negative_prompt, image=init_image).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-negative-1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative prompt = "ugly, deformed, disfigured, poor details, bad anatomy"</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative_prompt = "ugly, deformed, disfigured, poor details, bad anatomy"</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-negative-2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative prompt = "jungle"</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">negative_prompt = "jungle"</figcaption>
</div>
</div>
@@ -342,52 +331,54 @@ Start by generating an image with the text-to-image pipeline:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
from diffusers.utils import make_image_grid
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k").images[0]
text2image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k").images[0]
text2image
```
Now you can pass this generated image to the image-to-image pipeline:
```py
pipeline = AutoPipelineForImage2Image.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
"kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", image=image).images[0]
image
image2image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", image=text2image).images[0]
make_image_grid([text2image, image2image], rows=1, cols=2)
```
### Image-to-image-to-image
You can also chain multiple image-to-image pipelines together to create more interesting images. This can be useful for iteratively performing style transfer on an image, generate short GIFs, restore color to an image, or restore missing areas of an image.
You can also chain multiple image-to-image pipelines together to create more interesting images. This can be useful for iteratively performing style transfer on an image, generating short GIFs, restoring color to an image, or restoring missing areas of an image.
Start by generating an image:
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
@@ -404,10 +395,11 @@ It is important to specify `output_type="latent"` in the pipeline to keep all th
Pass the latent output from this pipeline to the next pipeline to generate an image in a [comic book art style](https://huggingface.co/ogkalu/Comic-Diffusion):
```py
pipelne = AutoPipelineForImage2Image.from_pretrained(
"ogkalu/Comic-Diffusion", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
pipeline = AutoPipelineForImage2Image.from_pretrained(
"ogkalu/Comic-Diffusion", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# need to include the token "charliebo artstyle" in the prompt to use this checkpoint
@@ -418,14 +410,15 @@ Repeat one more time to generate the final image in a [pixel art style](https://
```py
pipeline = AutoPipelineForImage2Image.from_pretrained(
"kohbanye/pixel-art-style", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
"kohbanye/pixel-art-style", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# need to include the token "pixelartstyle" in the prompt to use this checkpoint
image = pipeline("Astronaut in a jungle, pixelartstyle", image=image).images[0]
image
make_image_grid([init_image, image], rows=1, cols=2)
```
### Image-to-upscaler-to-super-resolution
@@ -436,21 +429,19 @@ Start with an image-to-image pipeline:
```py
import torch
import requests
from PIL import Image
from io import BytesIO
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import make_image_grid, load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
@@ -467,7 +458,9 @@ It is important to specify `output_type="latent"` in the pipeline to keep all th
Chain it to an upscaler pipeline to increase the image resolution:
```py
upscaler = AutoPipelineForImage2Image.from_pretrained(
from diffusers import StableDiffusionLatentUpscalePipeline
upscaler = StableDiffusionLatentUpscalePipeline.from_pretrained(
"stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
upscaler.enable_model_cpu_offload()
@@ -479,14 +472,16 @@ image_2 = upscaler(prompt, image=image_1, output_type="latent").images[0]
Finally, chain it to a super-resolution pipeline to further enhance the resolution:
```py
super_res = AutoPipelineForImage2Image.from_pretrained(
from diffusers import StableDiffusionUpscalePipeline
super_res = StableDiffusionUpscalePipeline.from_pretrained(
"stabilityai/stable-diffusion-x4-upscaler", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
super_res.enable_model_cpu_offload()
super_res.enable_xformers_memory_efficient_attention()
image_3 = upscaler(prompt, image=image_2).images[0]
image_3
image_3 = super_res(prompt, image=image_2).images[0]
make_image_grid([init_image, image_3.resize((512, 512))], rows=1, cols=2)
```
## Control image generation
@@ -504,13 +499,14 @@ from diffusers import AutoPipelineForImage2Image
import torch
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline(prompt_emebds=prompt_embeds, # generated from Compel
negative_prompt_embeds, # generated from Compel
image = pipeline(prompt_embeds=prompt_embeds, # generated from Compel
negative_prompt_embeds=negative_prompt_embeds, # generated from Compel
image=init_image,
).images[0]
```
@@ -522,19 +518,20 @@ ControlNets provide a more flexible and accurate way to control image generation
For example, let's condition an image with a depth map to keep the spatial information in the image.
```py
from diffusers.utils import load_image, make_image_grid
# prepare image
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png"
response = requests.get(url)
init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = load_image(url)
init_image = init_image.resize((958, 960)) # resize to depth image dimensions
depth_image = load_image("https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth/resolve/main/images/control.png")
make_image_grid([init_image, depth_image], rows=1, cols=2)
```
Load a ControlNet model conditioned on depth maps and the [`AutoPipelineForImage2Image`]:
```py
from diffusers import ControlNetModel, AutoPipelineForImage2Image
from diffusers.utils import load_image
import torch
controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11f1p_sd15_depth", torch_dtype=torch.float16, variant="fp16", use_safetensors=True)
@@ -542,6 +539,7 @@ pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
```
@@ -549,8 +547,8 @@ Now generate a new image conditioned on the depth map, initial image, and prompt
```py
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipeline(prompt, image=init_image, control_image=depth_image).images[0]
image
image_control_net = pipeline(prompt, image=init_image, control_image=depth_image).images[0]
make_image_grid([init_image, depth_image, image_control_net], rows=1, cols=3)
```
<div class="flex flex-row gap-4">
@@ -575,13 +573,14 @@ pipeline = AutoPipelineForImage2Image.from_pretrained(
"nitrosocke/elden-ring-diffusion", torch_dtype=torch.float16,
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
prompt = "elden ring style astronaut in a jungle" # include the token "elden ring style" in the prompt
negative_prompt = "ugly, deformed, disfigured, poor details, bad anatomy"
image = pipeline(prompt, negative_prompt=negative_prompt, image=init_image, strength=0.45, guidance_scale=10.5).images[0]
image
image_elden_ring = pipeline(prompt, negative_prompt=negative_prompt, image=image_control_net, strength=0.45, guidance_scale=10.5).images[0]
make_image_grid([init_image, depth_image, image_control_net, image_elden_ring], rows=2, cols=2)
```
<div class="flex justify-center">
@@ -597,10 +596,10 @@ Running diffusion models is computationally expensive and intensive, but with a
+ pipeline.enable_xformers_memory_efficient_attention()
```
With [`torch.compile`](../optimization/torch2.0#torch.compile), you can boost your inference speed even more by wrapping your UNet with it:
With [`torch.compile`](../optimization/torch2.0#torchcompile), you can boost your inference speed even more by wrapping your UNet with it:
```py
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipeline.unet = torch.compile(pipeline.unet, mode="reduce-overhead", fullgraph=True)
```
To learn more, take a look at the [Reduce memory usage](../optimization/memory) and [Torch 2.0](../optimization/torch2.0) guides.
+77 -49
View File
@@ -23,12 +23,13 @@ With 🤗 Diffusers, here is how you can do inpainting:
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder-inpaint", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
```
@@ -41,8 +42,8 @@ You'll notice throughout the guide, we use [`~DiffusionPipeline.enable_model_cpu
2. Load the base and mask images:
```py
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
```
3. Create a prompt to inpaint the image with and pass it to the pipeline with the base and mask images:
@@ -51,6 +52,7 @@ mask_image = load_image("https://huggingface.co/datasets/huggingface/documentati
prompt = "a black cat with glowing eyes, cute, adorable, disney, pixar, highly detailed, 8k"
negative_prompt = "bad anatomy, deformed, ugly, disfigured"
image = pipeline(prompt=prompt, negative_prompt=negative_prompt, image=init_image, mask_image=mask_image).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
<div class="flex gap-4">
@@ -58,6 +60,10 @@ image = pipeline(prompt=prompt, negative_prompt=negative_prompt, image=init_imag
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">base image</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">mask image</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint-cat.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
@@ -79,7 +85,7 @@ Upload a base image to inpaint on and use the sketch tool to draw a mask. Once y
## Popular models
[Stable Diffusion Inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
[Stable Diffusion Inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2 Inpainting](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
### Stable Diffusion Inpainting
@@ -88,21 +94,23 @@ Stable Diffusion Inpainting is a latent diffusion model finetuned on 512x512 ima
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
generator = torch.Generator("cuda").manual_seed(92)
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, generator=generator).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
### Stable Diffusion XL (SDXL) Inpainting
@@ -112,21 +120,23 @@ SDXL is a larger and more powerful version of Stable Diffusion v1.5. This model
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"diffusers/stable-diffusion-xl-1.0-inpainting-0.1", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
generator = torch.Generator("cuda").manual_seed(92)
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, generator=generator).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
### Kandinsky 2.2 Inpainting
@@ -136,21 +146,23 @@ The Kandinsky model family is similar to SDXL because it uses two models as well
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder-inpaint", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
generator = torch.Generator("cuda").manual_seed(92)
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, generator=generator).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
<div class="flex flex-row gap-4">
@@ -186,20 +198,22 @@ Image features - like quality and "creativity" - are dependent on pipeline param
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.6).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
<div class="flex flex-row gap-4">
@@ -229,20 +243,22 @@ You can use `strength` and `guidance_scale` together for more control over how e
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, guidance_scale=2.5).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
<div class="flex flex-row gap-4">
@@ -267,22 +283,23 @@ A negative prompt assumes the opposite role of a prompt; it guides the model awa
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
negative_prompt = "bad architecture, unstable, poor details, blurry"
image = pipeline(prompt=prompt, negative_prompt=negative_prompt, image=init_image, mask_image=mask_image).images[0]
image
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
<div class="flex justify-center">
@@ -302,7 +319,7 @@ import numpy as np
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
device = "cuda"
pipeline = AutoPipelineForInpainting.from_pretrained(
@@ -334,6 +351,7 @@ mask_image_arr[mask_image_arr >= 0.5] = 1
unmasked_unchanged_image_arr = (1 - mask_image_arr) * init_image + mask_image_arr * repainted_image
unmasked_unchanged_image = PIL.Image.fromarray(unmasked_unchanged_image_arr.round().astype("uint8"))
unmasked_unchanged_image.save("force_unmasked_unchanged.png")
make_image_grid([init_image, mask_image, repainted_image, unmasked_unchanged_image], rows=2, cols=2)
```
## Chained inpainting pipelines
@@ -349,35 +367,37 @@ Start with the text-to-image pipeline to create a castle:
```py
import torch
from diffusers import AutoPipelineForText2Image, AutoPipelineForInpainting
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline("concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k").images[0]
text2image = pipeline("concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k").images[0]
```
Load the mask image of the output from above:
```py
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_text-chain-mask.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_text-chain-mask.png")
```
And let's inpaint the masked area with a waterfall:
```py
pipeline = AutoPipelineForInpainting.from_pretrained(
"kandinsky-community/kandinsky-2-2-decoder-inpaint", torch_dtype=torch.float16, variant="fp16"
"kandinsky-community/kandinsky-2-2-decoder-inpaint", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
prompt = "digital painting of a fantasy waterfall, cloudy"
image = pipeline(prompt=prompt, image=image, mask_image=mask_image).images[0]
image
image = pipeline(prompt=prompt, image=text2image, mask_image=mask_image).images[0]
make_image_grid([text2image, mask_image, image], rows=1, cols=3)
```
<div class="flex flex-row gap-4">
@@ -391,7 +411,6 @@ image
</div>
</div>
### Inpaint-to-image-to-image
You can also chain an inpainting pipeline before another pipeline like image-to-image or an upscaler to improve the quality.
@@ -401,23 +420,24 @@ Begin by inpainting an image:
```py
import torch
from diffusers import AutoPipelineForInpainting, AutoPipelineForImage2Image
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
image_inpainting = pipeline(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
# resize image to 1024x1024 for SDXL
image = image.resize((1024, 1024))
image_inpainting = image_inpainting.resize((1024, 1024))
```
Now let's pass the image to another inpainting pipeline with SDXL's refiner model to enhance the image details and quality:
@@ -427,9 +447,10 @@ pipeline = AutoPipelineForInpainting.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline(prompt=prompt, image=image, mask_image=mask_image, output_type="latent").images[0]
image = pipeline(prompt=prompt, image=image_inpainting, mask_image=mask_image, output_type="latent").images[0]
```
<Tip>
@@ -442,9 +463,11 @@ Finally, you can pass this image to an image-to-image pipeline to put the finish
```py
pipeline = AutoPipelineForImage2Image.from_pipe(pipeline)
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline(prompt=prompt, image=image).images[0]
make_image_grid([init_image, mask_image, image_inpainting, image], rows=2, cols=2)
```
<div class="flex flex-row gap-4">
@@ -477,18 +500,21 @@ Once you've generated the embeddings, pass them to the `prompt_embeds` (and `neg
```py
import torch
from diffusers import AutoPipelineForInpainting
from diffusers.utils import make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16,
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
image = pipeline(prompt_emebds=prompt_embeds, # generated from Compel
negative_prompt_embeds, # generated from Compel
image = pipeline(prompt_embeds=prompt_embeds, # generated from Compel
negative_prompt_embeds=negative_prompt_embeds, # generated from Compel
image=init_image,
mask_image=mask_image
).images[0]
make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
### ControlNet
@@ -501,7 +527,7 @@ For example, let's condition an image with a ControlNet pretrained on inpaint im
import torch
import numpy as np
from diffusers import ControlNetModel, StableDiffusionControlNetInpaintPipeline
from diffusers.utils import load_image
from diffusers.utils import load_image, make_image_grid
# load ControlNet
controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_inpaint", torch_dtype=torch.float16, variant="fp16")
@@ -511,11 +537,12 @@ pipeline = StableDiffusionControlNetInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16"
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
# load base and mask image
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png").convert("RGB")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png").convert("RGB")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint_mask.png")
# prepare control image
def make_inpaint_condition(init_image, mask_image):
@@ -536,7 +563,7 @@ Now generate an image from the base, mask and control images. You'll notice feat
```py
prompt = "concept art digital painting of an elven castle, inspired by lord of the rings, highly detailed, 8k"
image = pipeline(prompt=prompt, image=init_image, mask_image=mask_image, control_image=control_image).images[0]
image
make_image_grid([init_image, mask_image, PIL.Image.fromarray(np.uint8(control_image[0][0])).convert('RGB'), image], rows=2, cols=2)
```
You can take this a step further and chain it with an image-to-image pipeline to apply a new [style](https://huggingface.co/nitrosocke/elden-ring-diffusion):
@@ -548,13 +575,14 @@ pipeline = AutoPipelineForImage2Image.from_pretrained(
"nitrosocke/elden-ring-diffusion", torch_dtype=torch.float16,
).to("cuda")
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
pipeline.enable_xformers_memory_efficient_attention()
prompt = "elden ring style castle" # include the token "elden ring style" in the prompt
negative_prompt = "bad architecture, deformed, disfigured, poor details"
image = pipeline(prompt, negative_prompt=negative_prompt, image=image).images[0]
image
image_elden_ring = pipeline(prompt, negative_prompt=negative_prompt, image=image).images[0]
make_image_grid([init_image, mask_image, image, image_elden_ring], rows=2, cols=2)
```
<div class="flex flex-row gap-4">
@@ -576,17 +604,17 @@ image
It can be difficult and slow to run diffusion models if you're resource constrained, but it doesn't have to be with a few optimization tricks. One of the biggest (and easiest) optimizations you can enable is switching to memory-efficient attention. If you're using PyTorch 2.0, [scaled-dot product attention](../optimization/torch2.0#scaled-dot-product-attention) is automatically enabled and you don't need to do anything else. For non-PyTorch 2.0 users, you can install and use [xFormers](../optimization/xformers)'s implementation of memory-efficient attention. Both options reduce memory usage and accelerate inference.
You can also offload the model to the GPU to save even more memory:
You can also offload the model to the CPU to save even more memory:
```diff
+ pipeline.enable_xformers_memory_efficient_attention()
+ pipeline.enable_model_cpu_offload()
```
To speed-up your inference code even more, use [`torch_compile`](../optimization/torch2.0#torch.compile). You should wrap `torch.compile` around the most intensive component in the pipeline which is typically the UNet:
To speed-up your inference code even more, use [`torch_compile`](../optimization/torch2.0#torchcompile). You should wrap `torch.compile` around the most intensive component in the pipeline which is typically the UNet:
```py
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipeline.unet = torch.compile(pipeline.unet, mode="reduce-overhead", fullgraph=True)
```
Learn more in the [Reduce memory usage](../optimization/memory) and [Torch 2.0](../optimization/torch2.0) guides.
Learn more in the [Reduce memory usage](../optimization/memory) and [Torch 2.0](../optimization/torch2.0) guides.
+154
View File
@@ -0,0 +1,154 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Performing inference with LCM
Latent Consistency Models (LCM) enable quality image generation in typically 2-4 steps making it possible to use diffusion models in almost real-time settings.
From the [official website](https://latent-consistency-models.github.io/):
> LCMs can be distilled from any pre-trained Stable Diffusion (SD) in only 4,000 training steps (~32 A100 GPU Hours) for generating high quality 768 x 768 resolution images in 2~4 steps or even one step, significantly accelerating text-to-image generation. We employ LCM to distill the Dreamshaper-V7 version of SD in just 4,000 training iterations.
For a more technical overview of LCMs, refer to [the paper](https://huggingface.co/papers/2310.04378).
This guide shows how to perform inference with LCMs for text-to-image and image-to-image generation tasks. It will also cover performing inference with LoRA checkpoints.
## Text-to-image
You'll use the [`StableDiffusionXLPipeline`] here changing the `unet`. The UNet was distilled from the SDXL UNet using the framework introduced in LCM. Another important component is the scheduler: [`LCMScheduler`]. Together with the distilled UNet and the scheduler, LCM enables a fast inference workflow overcoming the slow iterative nature of diffusion models.
```python
from diffusers import DiffusionPipeline, UNet2DConditionModel, LCMScheduler
import torch
unet = UNet2DConditionModel.from_pretrained(
"latent-consistency/lcm-sdxl",
torch_dtype=torch.float16,
variant="fp16",
)
pipe = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", unet=unet, torch_dtype=torch.float16
).to("cuda")
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
prompt = "Self-portrait oil painting, a beautiful cyborg with golden hair, 8k"
generator = torch.manual_seed(0)
image = pipe(
prompt=prompt, num_inference_steps=4, generator=generator, guidance_scale=8.0
).images[0]
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_intro.png)
Notice that we use only 4 steps for generation which is way less than what's typically used for standard SDXL.
Some details to keep in mind:
* To perform classifier-free guidance, batch size is usually doubled inside the pipeline. LCM, however, applies guidance using guidance embeddings, so the batch size does not have to be doubled in this case. This leads to a faster inference time, with the drawback that negative prompts don't have any effect on the denoising process.
* The UNet was trained using the [3., 13.] guidance scale range. So, that is the ideal range for `guidance_scale`. However, disabling `guidance_scale` using a value of 1.0 is also effective in most cases.
## Image-to-image
The findings above apply to image-to-image tasks too. Let's look at how we can perform image-to-image generation with LCMs:
```python
from diffusers import AutoPipelineForImage2Image, UNet2DConditionModel, LCMScheduler
from diffusers.utils import load_image
import torch
unet = UNet2DConditionModel.from_pretrained(
"latent-consistency/lcm-sdxl",
torch_dtype=torch.float16,
variant="fp16",
)
pipe = AutoPipelineForImage2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", unet=unet, torch_dtype=torch.float16
).to("cuda")
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
prompt = "High altitude snowy mountains"
image = load_image(
"https://huggingface.co/datasets/sayakpaul/sample-datasets/resolve/main/snowy_mountains.jpeg"
)
generator = torch.manual_seed(0)
image = pipe(
prompt=prompt,
image=image,
num_inference_steps=4,
generator=generator,
guidance_scale=8.0,
).images[0]
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_i2i.png)
## LoRA
It is possible to generalize the LCM framework to use with [LoRA](../training/lora.md). It effectively eliminates the need to conduct expensive fine-tuning runs as LoRA training concerns just a few number of parameters compared to full fine-tuning. During inference, the [`LCMScheduler`] comes to the advantage as it enables very few-steps inference without compromising the quality.
We recommend to disable `guidance_scale` by setting it 0. The model is trained to follow prompts accurately
even without using guidance scale. You can however, still use guidance scale in which case we recommend
using values between 1.0 and 2.0.
### Text-to-image
```python
from diffusers import DiffusionPipeline, LCMScheduler
import torch
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
lcm_lora_id = "latent-consistency/lcm-lora-sdxl"
pipe = DiffusionPipeline.from_pretrained(model_id, variant="fp16", torch_dtype=torch.float16).to("cuda")
pipe.load_lora_weights(lcm_lora_id)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
prompt = "close-up photography of old man standing in the rain at night, in a street lit by lamps, leica 35mm summilux"
image = pipe(
prompt=prompt,
num_inference_steps=4,
guidance_scale=0, # set guidance scale to 0 to disable it
).images[0]
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lora_lcm.png)
### Image-to-image
Extending LCM LoRA to image-to-image is possible:
```python
from diffusers import StableDiffusionXLImg2ImgPipeline, LCMScheduler
from diffusers.utils import load_image
import torch
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
lcm_lora_id = "latent-consistency/lcm-lora-sdxl"
pipe = StableDiffusionXLImg2ImgPipeline.from_pretrained(model_id, variant="fp16", torch_dtype=torch.float16).to("cuda")
pipe.load_lora_weights(lcm_lora_id)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
prompt = "close-up photography of old man standing in the rain at night, in a street lit by lamps, leica 35mm summilux"
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lora_lcm.png")
image = pipe(
prompt=prompt,
image=image,
num_inference_steps=4,
guidance_scale=0, # set guidance scale to 0 to disable it
).images[0]
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lcm/lcm_lora_i2i.png)
@@ -23,16 +23,16 @@ You can use any of the 🧨 Diffusers [checkpoints](https://huggingface.co/model
<Tip>
💡 Want to train your own unconditional image generation model? Take a look at the training [guide](training/unconditional_training) to learn how to generate your own images.
💡 Want to train your own unconditional image generation model? Take a look at the training [guide](../training/unconditional_training) to learn how to generate your own images.
</Tip>
In this guide, you'll use [`DiffusionPipeline`] for unconditional image generation with [DDPM](https://arxiv.org/abs/2006.11239):
```python
>>> from diffusers import DiffusionPipeline
from diffusers import DiffusionPipeline
>>> generator = DiffusionPipeline.from_pretrained("anton-l/ddpm-butterflies-128", use_safetensors=True)
generator = DiffusionPipeline.from_pretrained("anton-l/ddpm-butterflies-128", use_safetensors=True)
```
The [`DiffusionPipeline`] downloads and caches all modeling, tokenization, and scheduling components.
@@ -40,13 +40,14 @@ Because the model consists of roughly 1.4 billion parameters, we strongly recomm
You can move the generator object to a GPU, just like you would in PyTorch:
```python
>>> generator.to("cuda")
generator.to("cuda")
```
Now you can use the `generator` to generate an image:
```python
>>> image = generator().images[0]
image = generator().images[0]
image
```
The output is by default wrapped into a [`PIL.Image`](https://pillow.readthedocs.io/en/stable/reference/Image.html?highlight=image#the-image-class) object.
@@ -54,7 +55,7 @@ The output is by default wrapped into a [`PIL.Image`](https://pillow.readthedocs
You can save the image by calling:
```python
>>> image.save("generated_image.png")
image.save("generated_image.png")
```
Try out the Spaces below, and feel free to play around with the inference steps parameter to see how it affects the image quality!
@@ -65,5 +66,3 @@ Try out the Spaces below, and feel free to play around with the inference steps
width="850"
height="500"
></iframe>
+1 -1
View File
@@ -56,7 +56,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
+1 -1
View File
@@ -59,7 +59,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = logging.getLogger(__name__)
+1 -1
View File
@@ -58,7 +58,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -62,7 +62,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
+1 -1
View File
@@ -61,7 +61,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
+1 -1
View File
@@ -35,7 +35,7 @@ from diffusers.utils import check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
# Cache compiled models across invocations of this script.
cc.initialize_cache(os.path.expanduser("~/.cache/jax/compilation_cache"))
+1 -1
View File
@@ -68,7 +68,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -58,7 +58,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -52,7 +52,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -55,7 +55,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -52,7 +52,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -51,7 +51,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -58,7 +58,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -53,7 +53,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -33,7 +33,7 @@ from diffusers.utils import check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = logging.getLogger(__name__)
@@ -49,7 +49,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -58,7 +58,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -57,7 +57,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -79,7 +79,7 @@ else:
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__)
@@ -56,7 +56,7 @@ else:
# ------------------------------------------------------------------------------
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = logging.getLogger(__name__)
@@ -29,7 +29,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -50,7 +50,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
@@ -51,7 +51,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.23.0.dev0")
check_min_version("0.23.0")
logger = get_logger(__name__, log_level="INFO")
File diff suppressed because it is too large Load Diff
+3 -1
View File
@@ -112,10 +112,12 @@ _deps = [
"numpy",
"omegaconf",
"parameterized",
"peft<=0.6.2",
"protobuf>=3.20.3,<4",
"pytest",
"pytest-timeout",
"pytest-xdist",
"python>=3.8.0",
"ruff==0.0.280",
"safetensors>=0.3.1",
"sentencepiece>=0.1.91,!=0.1.92",
@@ -244,7 +246,7 @@ install_requires = [
setup(
name="diffusers",
version="0.23.0.dev0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
version="0.23.1", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
description="State-of-the-art diffusion in PyTorch and JAX.",
long_description=open("README.md", "r", encoding="utf-8").read(),
long_description_content_type="text/markdown",
+3 -1
View File
@@ -1,4 +1,4 @@
__version__ = "0.23.0.dev0"
__version__ = "0.23.1"
from typing import TYPE_CHECKING
@@ -77,6 +77,7 @@ else:
"AsymmetricAutoencoderKL",
"AutoencoderKL",
"AutoencoderTiny",
"ConsistencyDecoderVAE",
"ControlNetModel",
"ModelMixin",
"MotionAdapter",
@@ -443,6 +444,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
AsymmetricAutoencoderKL,
AutoencoderKL,
AutoencoderTiny,
ConsistencyDecoderVAE,
ControlNetModel,
ModelMixin,
MotionAdapter,
+1 -13
View File
@@ -23,21 +23,9 @@ from .utils.versions import require_version, require_version_core
# order specific notes:
# - tqdm must be checked before tokenizers
pkgs_to_check_at_runtime = "python tqdm regex requests packaging filelock numpy tokenizers".split()
if sys.version_info < (3, 7):
pkgs_to_check_at_runtime.append("dataclasses")
if sys.version_info < (3, 8):
pkgs_to_check_at_runtime.append("importlib_metadata")
pkgs_to_check_at_runtime = "python requests filelock numpy".split()
for pkg in pkgs_to_check_at_runtime:
if pkg in deps:
if pkg == "tokenizers":
# must be loaded here, or else tqdm check may fail
from .utils import is_tokenizers_available
if not is_tokenizers_available():
continue # not required, check version only if installed
require_version_core(deps[pkg])
else:
raise ValueError(f"can't find {pkg} in {deps.keys()}, check dependency_versions_table.py")
@@ -25,10 +25,12 @@ deps = {
"numpy": "numpy",
"omegaconf": "omegaconf",
"parameterized": "parameterized",
"peft": "peft<=0.6.2",
"protobuf": "protobuf>=3.20.3,<4",
"pytest": "pytest",
"pytest-timeout": "pytest-timeout",
"pytest-xdist": "pytest-xdist",
"python": "python>=3.8.0",
"ruff": "ruff==0.0.280",
"safetensors": "safetensors>=0.3.1",
"sentencepiece": "sentencepiece>=0.1.91,!=0.1.92",
+8 -3
View File
@@ -1411,6 +1411,11 @@ class LoraLoaderMixin:
filter(lambda x: all(substring not in x for substring in unallowed_substrings), targeted_files)
)
if any(f.endswith(LORA_WEIGHT_NAME) for f in targeted_files):
targeted_files = list(filter(lambda x: x.endswith(LORA_WEIGHT_NAME), targeted_files))
elif any(f.endswith(LORA_WEIGHT_NAME_SAFE) for f in targeted_files):
targeted_files = list(filter(lambda x: x.endswith(LORA_WEIGHT_NAME_SAFE), targeted_files))
if len(targeted_files) > 1:
raise ValueError(
f"Provided path contains more than one weights file in the {file_extension} format. Either specify `weight_name` in `load_lora_weights` or make sure there's only one `.safetensors` or `.bin` file in {pretrained_model_name_or_path_or_dict}."
@@ -2390,7 +2395,7 @@ class LoraLoaderMixin:
def set_adapters_for_text_encoder(
self,
adapter_names: Union[List[str], str],
text_encoder: Optional[PreTrainedModel] = None,
text_encoder: Optional["PreTrainedModel"] = None, # noqa: F821
text_encoder_weights: List[float] = None,
):
"""
@@ -2429,7 +2434,7 @@ class LoraLoaderMixin:
)
set_weights_and_activate_adapters(text_encoder, adapter_names, text_encoder_weights)
def disable_lora_for_text_encoder(self, text_encoder: Optional[PreTrainedModel] = None):
def disable_lora_for_text_encoder(self, text_encoder: Optional["PreTrainedModel"] = None):
"""
Disables the LoRA layers for the text encoder.
@@ -2446,7 +2451,7 @@ class LoraLoaderMixin:
raise ValueError("Text Encoder not found.")
set_adapter_layers(text_encoder, enabled=False)
def enable_lora_for_text_encoder(self, text_encoder: Optional[PreTrainedModel] = None):
def enable_lora_for_text_encoder(self, text_encoder: Optional["PreTrainedModel"] = None):
"""
Enables the LoRA layers for the text encoder.
+2
View File
@@ -24,6 +24,7 @@ if is_torch_available():
_import_structure["autoencoder_asym_kl"] = ["AsymmetricAutoencoderKL"]
_import_structure["autoencoder_kl"] = ["AutoencoderKL"]
_import_structure["autoencoder_tiny"] = ["AutoencoderTiny"]
_import_structure["consistency_decoder_vae"] = ["ConsistencyDecoderVAE"]
_import_structure["controlnet"] = ["ControlNetModel"]
_import_structure["dual_transformer_2d"] = ["DualTransformer2DModel"]
_import_structure["modeling_utils"] = ["ModelMixin"]
@@ -50,6 +51,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
from .autoencoder_asym_kl import AsymmetricAutoencoderKL
from .autoencoder_kl import AutoencoderKL
from .autoencoder_tiny import AutoencoderTiny
from .consistency_decoder_vae import ConsistencyDecoderVAE
from .controlnet import ControlNetModel
from .dual_transformer_2d import DualTransformer2DModel
from .modeling_utils import ModelMixin
+1 -1
View File
@@ -287,7 +287,7 @@ class BasicTransformerBlock(nn.Module):
else:
raise ValueError("Incorrect norm")
if self.pos_embed is not None and self.use_ada_layer_norm_single is None:
if self.pos_embed is not None and self.use_ada_layer_norm_single is False:
norm_hidden_states = self.pos_embed(norm_hidden_states)
attn_output = self.attn2(
+28 -22
View File
@@ -378,7 +378,7 @@ class Attention(nn.Module):
_remove_lora (`bool`, *optional*, defaults to `False`):
Set to `True` to remove LoRA layers from the model.
"""
if hasattr(self, "processor") and _remove_lora and self.to_q.lora_layer is not None:
if not USE_PEFT_BACKEND and hasattr(self, "processor") and _remove_lora and self.to_q.lora_layer is not None:
deprecate(
"set_processor to offload LoRA",
"0.26.0",
@@ -879,6 +879,9 @@ class AttnAddedKVProcessor:
scale: float = 1.0,
) -> torch.Tensor:
residual = hidden_states
args = () if USE_PEFT_BACKEND else (scale,)
hidden_states = hidden_states.view(hidden_states.shape[0], hidden_states.shape[1], -1).transpose(1, 2)
batch_size, sequence_length, _ = hidden_states.shape
@@ -891,17 +894,17 @@ class AttnAddedKVProcessor:
hidden_states = attn.group_norm(hidden_states.transpose(1, 2)).transpose(1, 2)
query = attn.to_q(hidden_states, scale=scale)
query = attn.to_q(hidden_states, *args)
query = attn.head_to_batch_dim(query)
encoder_hidden_states_key_proj = attn.add_k_proj(encoder_hidden_states, scale=scale)
encoder_hidden_states_value_proj = attn.add_v_proj(encoder_hidden_states, scale=scale)
encoder_hidden_states_key_proj = attn.add_k_proj(encoder_hidden_states, *args)
encoder_hidden_states_value_proj = attn.add_v_proj(encoder_hidden_states, *args)
encoder_hidden_states_key_proj = attn.head_to_batch_dim(encoder_hidden_states_key_proj)
encoder_hidden_states_value_proj = attn.head_to_batch_dim(encoder_hidden_states_value_proj)
if not attn.only_cross_attention:
key = attn.to_k(hidden_states, scale=scale)
value = attn.to_v(hidden_states, scale=scale)
key = attn.to_k(hidden_states, *args)
value = attn.to_v(hidden_states, *args)
key = attn.head_to_batch_dim(key)
value = attn.head_to_batch_dim(value)
key = torch.cat([encoder_hidden_states_key_proj, key], dim=1)
@@ -915,7 +918,7 @@ class AttnAddedKVProcessor:
hidden_states = attn.batch_to_head_dim(hidden_states)
# linear proj
hidden_states = attn.to_out[0](hidden_states, scale=scale)
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
@@ -946,6 +949,9 @@ class AttnAddedKVProcessor2_0:
scale: float = 1.0,
) -> torch.Tensor:
residual = hidden_states
args = () if USE_PEFT_BACKEND else (scale,)
hidden_states = hidden_states.view(hidden_states.shape[0], hidden_states.shape[1], -1).transpose(1, 2)
batch_size, sequence_length, _ = hidden_states.shape
@@ -958,7 +964,7 @@ class AttnAddedKVProcessor2_0:
hidden_states = attn.group_norm(hidden_states.transpose(1, 2)).transpose(1, 2)
query = attn.to_q(hidden_states, scale=scale)
query = attn.to_q(hidden_states, *args)
query = attn.head_to_batch_dim(query, out_dim=4)
encoder_hidden_states_key_proj = attn.add_k_proj(encoder_hidden_states)
@@ -967,8 +973,8 @@ class AttnAddedKVProcessor2_0:
encoder_hidden_states_value_proj = attn.head_to_batch_dim(encoder_hidden_states_value_proj, out_dim=4)
if not attn.only_cross_attention:
key = attn.to_k(hidden_states, scale=scale)
value = attn.to_v(hidden_states, scale=scale)
key = attn.to_k(hidden_states, *args)
value = attn.to_v(hidden_states, *args)
key = attn.head_to_batch_dim(key, out_dim=4)
value = attn.head_to_batch_dim(value, out_dim=4)
key = torch.cat([encoder_hidden_states_key_proj, key], dim=2)
@@ -985,7 +991,7 @@ class AttnAddedKVProcessor2_0:
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, residual.shape[1])
# linear proj
hidden_states = attn.to_out[0](hidden_states, scale=scale)
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
@@ -1177,6 +1183,8 @@ class AttnProcessor2_0:
) -> torch.FloatTensor:
residual = hidden_states
args = () if USE_PEFT_BACKEND else (scale,)
if attn.spatial_norm is not None:
hidden_states = attn.spatial_norm(hidden_states, temb)
@@ -1207,12 +1215,8 @@ class AttnProcessor2_0:
elif attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
key = (
attn.to_k(encoder_hidden_states, scale=scale) if not USE_PEFT_BACKEND else attn.to_k(encoder_hidden_states)
)
value = (
attn.to_v(encoder_hidden_states, scale=scale) if not USE_PEFT_BACKEND else attn.to_v(encoder_hidden_states)
)
key = attn.to_k(encoder_hidden_states, *args)
value = attn.to_v(encoder_hidden_states, *args)
inner_dim = key.shape[-1]
head_dim = inner_dim // attn.heads
@@ -1232,9 +1236,7 @@ class AttnProcessor2_0:
hidden_states = hidden_states.to(query.dtype)
# linear proj
hidden_states = (
attn.to_out[0](hidden_states, scale=scale) if not USE_PEFT_BACKEND else attn.to_out[0](hidden_states)
)
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
@@ -1361,6 +1363,7 @@ class CustomDiffusionXFormersAttnProcessor(nn.Module):
hidden_states = attn.to_out[0](hidden_states)
# dropout
hidden_states = attn.to_out[1](hidden_states)
return hidden_states
@@ -1433,8 +1436,11 @@ class CustomDiffusionAttnProcessor2_0(nn.Module):
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
if self.train_kv:
key = self.to_k_custom_diffusion(encoder_hidden_states)
value = self.to_v_custom_diffusion(encoder_hidden_states)
key = self.to_k_custom_diffusion(encoder_hidden_states.to(self.to_k_custom_diffusion.weight.dtype))
value = self.to_v_custom_diffusion(encoder_hidden_states.to(self.to_v_custom_diffusion.weight.dtype))
key = key.to(attn.to_q.weight.dtype)
value = value.to(attn.to_q.weight.dtype)
else:
key = attn.to_k(encoder_hidden_states)
value = attn.to_v(encoder_hidden_states)
@@ -138,6 +138,7 @@ class AsymmetricAutoencoderKL(ModelMixin, ConfigMixin):
def decode(
self,
z: torch.FloatTensor,
generator: Optional[torch.Generator] = None,
image: Optional[torch.FloatTensor] = None,
mask: Optional[torch.FloatTensor] = None,
return_dict: bool = True,
+3 -1
View File
@@ -294,7 +294,9 @@ class AutoencoderKL(ModelMixin, ConfigMixin, FromOriginalVAEMixin):
return DecoderOutput(sample=dec)
@apply_forward_hook
def decode(self, z: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, torch.FloatTensor]:
def decode(
self, z: torch.FloatTensor, return_dict: bool = True, generator=None
) -> Union[DecoderOutput, torch.FloatTensor]:
"""
Decode a batch of images.
+4 -2
View File
@@ -14,7 +14,7 @@
from dataclasses import dataclass
from typing import Tuple, Union
from typing import Optional, Tuple, Union
import torch
@@ -307,7 +307,9 @@ class AutoencoderTiny(ModelMixin, ConfigMixin):
return AutoencoderTinyOutput(latents=output)
@apply_forward_hook
def decode(self, x: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, Tuple[torch.FloatTensor]]:
def decode(
self, x: torch.FloatTensor, generator: Optional[torch.Generator] = None, return_dict: bool = True
) -> Union[DecoderOutput, Tuple[torch.FloatTensor]]:
if self.use_slicing and x.shape[0] > 1:
output = [self._tiled_decode(x_slice) if self.use_tiling else self.decoder(x) for x_slice in x.split(1)]
output = torch.cat(output)
@@ -0,0 +1,430 @@
# Copyright 2023 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
from dataclasses import dataclass
from typing import Dict, Optional, Tuple, Union
import torch
import torch.nn.functional as F
from torch import nn
from ..configuration_utils import ConfigMixin, register_to_config
from ..schedulers import ConsistencyDecoderScheduler
from ..utils import BaseOutput
from ..utils.accelerate_utils import apply_forward_hook
from ..utils.torch_utils import randn_tensor
from .attention_processor import (
ADDED_KV_ATTENTION_PROCESSORS,
CROSS_ATTENTION_PROCESSORS,
AttentionProcessor,
AttnAddedKVProcessor,
AttnProcessor,
)
from .modeling_utils import ModelMixin
from .unet_2d import UNet2DModel
from .vae import DecoderOutput, DiagonalGaussianDistribution, Encoder
@dataclass
class ConsistencyDecoderVAEOutput(BaseOutput):
"""
Output of encoding method.
Args:
latent_dist (`DiagonalGaussianDistribution`):
Encoded outputs of `Encoder` represented as the mean and logvar of `DiagonalGaussianDistribution`.
`DiagonalGaussianDistribution` allows for sampling latents from the distribution.
"""
latent_dist: "DiagonalGaussianDistribution"
class ConsistencyDecoderVAE(ModelMixin, ConfigMixin):
r"""
The consistency decoder used with DALL-E 3.
Examples:
```py
>>> import torch
>>> from diffusers import DiffusionPipeline, ConsistencyDecoderVAE
>>> vae = ConsistencyDecoderVAE.from_pretrained("openai/consistency-decoder", torch_dtype=pipe.torch_dtype)
>>> pipe = StableDiffusionPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", vae=vae, torch_dtype=torch.float16
... ).to("cuda")
>>> pipe("horse", generator=torch.manual_seed(0)).images
```
"""
@register_to_config
def __init__(
self,
scaling_factor=0.18215,
latent_channels=4,
encoder_act_fn="silu",
encoder_block_out_channels=(128, 256, 512, 512),
encoder_double_z=True,
encoder_down_block_types=(
"DownEncoderBlock2D",
"DownEncoderBlock2D",
"DownEncoderBlock2D",
"DownEncoderBlock2D",
),
encoder_in_channels=3,
encoder_layers_per_block=2,
encoder_norm_num_groups=32,
encoder_out_channels=4,
decoder_add_attention=False,
decoder_block_out_channels=(320, 640, 1024, 1024),
decoder_down_block_types=(
"ResnetDownsampleBlock2D",
"ResnetDownsampleBlock2D",
"ResnetDownsampleBlock2D",
"ResnetDownsampleBlock2D",
),
decoder_downsample_padding=1,
decoder_in_channels=7,
decoder_layers_per_block=3,
decoder_norm_eps=1e-05,
decoder_norm_num_groups=32,
decoder_num_train_timesteps=1024,
decoder_out_channels=6,
decoder_resnet_time_scale_shift="scale_shift",
decoder_time_embedding_type="learned",
decoder_up_block_types=(
"ResnetUpsampleBlock2D",
"ResnetUpsampleBlock2D",
"ResnetUpsampleBlock2D",
"ResnetUpsampleBlock2D",
),
):
super().__init__()
self.encoder = Encoder(
act_fn=encoder_act_fn,
block_out_channels=encoder_block_out_channels,
double_z=encoder_double_z,
down_block_types=encoder_down_block_types,
in_channels=encoder_in_channels,
layers_per_block=encoder_layers_per_block,
norm_num_groups=encoder_norm_num_groups,
out_channels=encoder_out_channels,
)
self.decoder_unet = UNet2DModel(
add_attention=decoder_add_attention,
block_out_channels=decoder_block_out_channels,
down_block_types=decoder_down_block_types,
downsample_padding=decoder_downsample_padding,
in_channels=decoder_in_channels,
layers_per_block=decoder_layers_per_block,
norm_eps=decoder_norm_eps,
norm_num_groups=decoder_norm_num_groups,
num_train_timesteps=decoder_num_train_timesteps,
out_channels=decoder_out_channels,
resnet_time_scale_shift=decoder_resnet_time_scale_shift,
time_embedding_type=decoder_time_embedding_type,
up_block_types=decoder_up_block_types,
)
self.decoder_scheduler = ConsistencyDecoderScheduler()
self.register_to_config(block_out_channels=encoder_block_out_channels)
self.register_buffer(
"means",
torch.tensor([0.38862467, 0.02253063, 0.07381133, -0.0171294])[None, :, None, None],
persistent=False,
)
self.register_buffer(
"stds", torch.tensor([0.9654121, 1.0440036, 0.76147926, 0.77022034])[None, :, None, None], persistent=False
)
self.quant_conv = nn.Conv2d(2 * latent_channels, 2 * latent_channels, 1)
self.use_slicing = False
self.use_tiling = False
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.enable_tiling
def enable_tiling(self, use_tiling: bool = True):
r"""
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow
processing larger images.
"""
self.use_tiling = use_tiling
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.disable_tiling
def disable_tiling(self):
r"""
Disable tiled VAE decoding. If `enable_tiling` was previously enabled, this method will go back to computing
decoding in one step.
"""
self.enable_tiling(False)
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.enable_slicing
def enable_slicing(self):
r"""
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
"""
self.use_slicing = True
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.disable_slicing
def disable_slicing(self):
r"""
Disable sliced VAE decoding. If `enable_slicing` was previously enabled, this method will go back to computing
decoding in one step.
"""
self.use_slicing = False
@property
# Copied from diffusers.models.unet_2d_condition.UNet2DConditionModel.attn_processors
def attn_processors(self) -> Dict[str, AttentionProcessor]:
r"""
Returns:
`dict` of attention processors: A dictionary containing all attention processors used in the model with
indexed by its weight name.
"""
# set recursively
processors = {}
def fn_recursive_add_processors(name: str, module: torch.nn.Module, processors: Dict[str, AttentionProcessor]):
if hasattr(module, "get_processor"):
processors[f"{name}.processor"] = module.get_processor(return_deprecated_lora=True)
for sub_name, child in module.named_children():
fn_recursive_add_processors(f"{name}.{sub_name}", child, processors)
return processors
for name, module in self.named_children():
fn_recursive_add_processors(name, module, processors)
return processors
# Copied from diffusers.models.unet_2d_condition.UNet2DConditionModel.set_attn_processor
def set_attn_processor(
self, processor: Union[AttentionProcessor, Dict[str, AttentionProcessor]], _remove_lora=False
):
r"""
Sets the attention processor to use to compute attention.
Parameters:
processor (`dict` of `AttentionProcessor` or only `AttentionProcessor`):
The instantiated processor class or a dictionary of processor classes that will be set as the processor
for **all** `Attention` layers.
If `processor` is a dict, the key needs to define the path to the corresponding cross attention
processor. This is strongly recommended when setting trainable attention processors.
"""
count = len(self.attn_processors.keys())
if isinstance(processor, dict) and len(processor) != count:
raise ValueError(
f"A dict of processors was passed, but the number of processors {len(processor)} does not match the"
f" number of attention layers: {count}. Please make sure to pass {count} processor classes."
)
def fn_recursive_attn_processor(name: str, module: torch.nn.Module, processor):
if hasattr(module, "set_processor"):
if not isinstance(processor, dict):
module.set_processor(processor, _remove_lora=_remove_lora)
else:
module.set_processor(processor.pop(f"{name}.processor"), _remove_lora=_remove_lora)
for sub_name, child in module.named_children():
fn_recursive_attn_processor(f"{name}.{sub_name}", child, processor)
for name, module in self.named_children():
fn_recursive_attn_processor(name, module, processor)
# Copied from diffusers.models.unet_2d_condition.UNet2DConditionModel.set_default_attn_processor
def set_default_attn_processor(self):
"""
Disables custom attention processors and sets the default attention implementation.
"""
if all(proc.__class__ in ADDED_KV_ATTENTION_PROCESSORS for proc in self.attn_processors.values()):
processor = AttnAddedKVProcessor()
elif all(proc.__class__ in CROSS_ATTENTION_PROCESSORS for proc in self.attn_processors.values()):
processor = AttnProcessor()
else:
raise ValueError(
f"Cannot call `set_default_attn_processor` when attention processors are of type {next(iter(self.attn_processors.values()))}"
)
self.set_attn_processor(processor, _remove_lora=True)
@apply_forward_hook
def encode(
self, x: torch.FloatTensor, return_dict: bool = True
) -> Union[ConsistencyDecoderVAEOutput, Tuple[DiagonalGaussianDistribution]]:
"""
Encode a batch of images into latents.
Args:
x (`torch.FloatTensor`): Input batch of images.
return_dict (`bool`, *optional*, defaults to `True`):
Whether to return a [`~models.consistecy_decoder_vae.ConsistencyDecoderOoutput`] instead of a plain
tuple.
Returns:
The latent representations of the encoded images. If `return_dict` is True, a
[`~models.consistency_decoder_vae.ConsistencyDecoderVAEOutput`] is returned, otherwise a plain `tuple`
is returned.
"""
if self.use_tiling and (x.shape[-1] > self.tile_sample_min_size or x.shape[-2] > self.tile_sample_min_size):
return self.tiled_encode(x, return_dict=return_dict)
if self.use_slicing and x.shape[0] > 1:
encoded_slices = [self.encoder(x_slice) for x_slice in x.split(1)]
h = torch.cat(encoded_slices)
else:
h = self.encoder(x)
moments = self.quant_conv(h)
posterior = DiagonalGaussianDistribution(moments)
if not return_dict:
return (posterior,)
return ConsistencyDecoderVAEOutput(latent_dist=posterior)
@apply_forward_hook
def decode(
self,
z: torch.FloatTensor,
generator: Optional[torch.Generator] = None,
return_dict: bool = True,
num_inference_steps=2,
) -> Union[DecoderOutput, torch.FloatTensor]:
z = (z * self.config.scaling_factor - self.means) / self.stds
scale_factor = 2 ** (len(self.config.block_out_channels) - 1)
z = F.interpolate(z, mode="nearest", scale_factor=scale_factor)
batch_size, _, height, width = z.shape
self.decoder_scheduler.set_timesteps(num_inference_steps, device=self.device)
x_t = self.decoder_scheduler.init_noise_sigma * randn_tensor(
(batch_size, 3, height, width), generator=generator, dtype=z.dtype, device=z.device
)
for t in self.decoder_scheduler.timesteps:
model_input = torch.concat([self.decoder_scheduler.scale_model_input(x_t, t), z], dim=1)
model_output = self.decoder_unet(model_input, t).sample[:, :3, :, :]
prev_sample = self.decoder_scheduler.step(model_output, t, x_t, generator).prev_sample
x_t = prev_sample
x_0 = x_t
if not return_dict:
return (x_0,)
return DecoderOutput(sample=x_0)
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.blend_v
def blend_v(self, a, b, blend_extent):
blend_extent = min(a.shape[2], b.shape[2], blend_extent)
for y in range(blend_extent):
b[:, :, y, :] = a[:, :, -blend_extent + y, :] * (1 - y / blend_extent) + b[:, :, y, :] * (y / blend_extent)
return b
# Copied from diffusers.models.autoencoder_kl.AutoencoderKL.blend_h
def blend_h(self, a, b, blend_extent):
blend_extent = min(a.shape[3], b.shape[3], blend_extent)
for x in range(blend_extent):
b[:, :, :, x] = a[:, :, :, -blend_extent + x] * (1 - x / blend_extent) + b[:, :, :, x] * (x / blend_extent)
return b
def tiled_encode(self, x: torch.FloatTensor, return_dict: bool = True) -> ConsistencyDecoderVAEOutput:
r"""Encode a batch of images using a tiled encoder.
When this option is enabled, the VAE will split the input tensor into tiles to compute encoding in several
steps. This is useful to keep memory use constant regardless of image size. The end result of tiled encoding is
different from non-tiled encoding because each tile uses a different encoder. To avoid tiling artifacts, the
tiles overlap and are blended together to form a smooth output. You may still see tile-sized changes in the
output, but they should be much less noticeable.
Args:
x (`torch.FloatTensor`): Input batch of images.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~models.consistency_decoder_vae.ConsistencyDecoderVAEOutput`] instead of a
plain tuple.
Returns:
[`~models.consistency_decoder_vae.ConsistencyDecoderVAEOutput`] or `tuple`:
If return_dict is True, a [`~models.consistency_decoder_vae.ConsistencyDecoderVAEOutput`] is returned,
otherwise a plain `tuple` is returned.
"""
overlap_size = int(self.tile_sample_min_size * (1 - self.tile_overlap_factor))
blend_extent = int(self.tile_latent_min_size * self.tile_overlap_factor)
row_limit = self.tile_latent_min_size - blend_extent
# Split the image into 512x512 tiles and encode them separately.
rows = []
for i in range(0, x.shape[2], overlap_size):
row = []
for j in range(0, x.shape[3], overlap_size):
tile = x[:, :, i : i + self.tile_sample_min_size, j : j + self.tile_sample_min_size]
tile = self.encoder(tile)
tile = self.quant_conv(tile)
row.append(tile)
rows.append(row)
result_rows = []
for i, row in enumerate(rows):
result_row = []
for j, tile in enumerate(row):
# blend the above tile and the left tile
# to the current tile and add the current tile to the result row
if i > 0:
tile = self.blend_v(rows[i - 1][j], tile, blend_extent)
if j > 0:
tile = self.blend_h(row[j - 1], tile, blend_extent)
result_row.append(tile[:, :, :row_limit, :row_limit])
result_rows.append(torch.cat(result_row, dim=3))
moments = torch.cat(result_rows, dim=2)
posterior = DiagonalGaussianDistribution(moments)
if not return_dict:
return (posterior,)
return ConsistencyDecoderVAEOutput(latent_dist=posterior)
def forward(
self,
sample: torch.FloatTensor,
sample_posterior: bool = False,
return_dict: bool = True,
generator: Optional[torch.Generator] = None,
) -> Union[DecoderOutput, torch.FloatTensor]:
r"""
Args:
sample (`torch.FloatTensor`): Input sample.
sample_posterior (`bool`, *optional*, defaults to `False`):
Whether to sample from the posterior.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`DecoderOutput`] instead of a plain tuple.
"""
x = sample
posterior = self.encode(x).latent_dist
if sample_posterior:
z = posterior.sample(generator=generator)
else:
z = posterior.mode()
dec = self.decode(z, generator=generator).sample
if not return_dict:
return (dec,)
return DecoderOutput(sample=dec)
+16 -4
View File
@@ -778,16 +778,22 @@ class Conv1dBlock(nn.Module):
out_channels (`int`): Number of output channels.
kernel_size (`int` or `tuple`): Size of the convolving kernel.
n_groups (`int`, default `8`): Number of groups to separate the channels into.
activation (`str`, defaults `mish`): Name of the activation function.
"""
def __init__(
self, inp_channels: int, out_channels: int, kernel_size: Union[int, Tuple[int, int]], n_groups: int = 8
self,
inp_channels: int,
out_channels: int,
kernel_size: Union[int, Tuple[int, int]],
n_groups: int = 8,
activation: str = "mish",
):
super().__init__()
self.conv1d = nn.Conv1d(inp_channels, out_channels, kernel_size, padding=kernel_size // 2)
self.group_norm = nn.GroupNorm(n_groups, out_channels)
self.mish = nn.Mish()
self.mish = get_activation(activation)
def forward(self, inputs: torch.Tensor) -> torch.Tensor:
intermediate_repr = self.conv1d(inputs)
@@ -808,16 +814,22 @@ class ResidualTemporalBlock1D(nn.Module):
out_channels (`int`): Number of output channels.
embed_dim (`int`): Embedding dimension.
kernel_size (`int` or `tuple`): Size of the convolving kernel.
activation (`str`, defaults `mish`): It is possible to choose the right activation function.
"""
def __init__(
self, inp_channels: int, out_channels: int, embed_dim: int, kernel_size: Union[int, Tuple[int, int]] = 5
self,
inp_channels: int,
out_channels: int,
embed_dim: int,
kernel_size: Union[int, Tuple[int, int]] = 5,
activation: str = "mish",
):
super().__init__()
self.conv_in = Conv1dBlock(inp_channels, out_channels, kernel_size)
self.conv_out = Conv1dBlock(out_channels, out_channels, kernel_size)
self.time_emb_act = nn.Mish()
self.time_emb_act = get_activation(activation)
self.time_emb = nn.Linear(embed_dim, out_channels)
self.residual_conv = (
+3 -1
View File
@@ -339,6 +339,7 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
elif self.is_input_vectorized:
hidden_states = self.latent_image_embedding(hidden_states)
elif self.is_input_patches:
height, width = hidden_states.shape[-2] // self.patch_size, hidden_states.shape[-1] // self.patch_size
hidden_states = self.pos_embed(hidden_states)
if self.adaln_single is not None:
@@ -425,7 +426,8 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
hidden_states = hidden_states.squeeze(1)
# unpatchify
height = width = int(hidden_states.shape[1] ** 0.5)
if self.adaln_single is None:
height = width = int(hidden_states.shape[1] ** 0.5)
hidden_states = hidden_states.reshape(
shape=(-1, height, width, self.patch_size, self.patch_size, self.out_channels)
)
+4
View File
@@ -117,6 +117,7 @@ class UNet2DModel(ModelMixin, ConfigMixin):
add_attention: bool = True,
class_embed_type: Optional[str] = None,
num_class_embeds: Optional[int] = None,
num_train_timesteps: Optional[int] = None,
):
super().__init__()
@@ -144,6 +145,9 @@ class UNet2DModel(ModelMixin, ConfigMixin):
elif time_embedding_type == "positional":
self.time_proj = Timesteps(block_out_channels[0], flip_sin_to_cos, freq_shift)
timestep_input_dim = block_out_channels[0]
elif time_embedding_type == "learned":
self.time_proj = nn.Embedding(num_train_timesteps, block_out_channels[0])
timestep_input_dim = block_out_channels[0]
self.time_embedding = TimestepEmbedding(timestep_input_dim, time_embed_dim)
+2 -2
View File
@@ -162,8 +162,8 @@ class VQModel(ModelMixin, ConfigMixin):
If return_dict is True, a [`~models.vq_model.VQEncoderOutput`] is returned, otherwise a plain `tuple`
is returned.
"""
x = sample
h = self.encode(x).latents
h = self.encode(sample).latents
dec = self.decode(h).sample
if not return_dict:
@@ -588,6 +588,34 @@ class AltDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -605,7 +633,7 @@ class AltDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -804,6 +832,14 @@ class AltDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 6.5 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
# 7. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
@@ -818,6 +854,7 @@ class AltDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
return_dict=False,
)[0]
@@ -852,7 +889,9 @@ class AltDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
callback(step_idx, t, latents)
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -646,6 +646,34 @@ class AltDiffusionImg2ImgPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -659,7 +687,7 @@ class AltDiffusionImg2ImgPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -849,6 +877,14 @@ class AltDiffusionImg2ImgPipeline(
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 7.5 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
# 8. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
@@ -863,6 +899,7 @@ class AltDiffusionImg2ImgPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
return_dict=False,
)[0]
@@ -893,7 +930,9 @@ class AltDiffusionImg2ImgPipeline(
callback(step_idx, t, latents)
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -498,7 +498,7 @@ class AnimateDiffPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraLo
@torch.no_grad()
def __call__(
self,
prompt: Union[str, List[str]],
prompt: Union[str, List[str]] = None,
num_frames: Optional[int] = 16,
height: Optional[int] = None,
width: Optional[int] = None,
@@ -6,7 +6,9 @@ from ...utils import (
)
_import_structure = {"pipeline_consistency_models": ["ConsistencyModelPipeline"]}
_import_structure = {
"pipeline_consistency_models": ["ConsistencyModelPipeline"],
}
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
from .pipeline_consistency_models import ConsistencyModelPipeline
@@ -1,3 +1,17 @@
# Copyright 2023 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
from typing import Callable, List, Optional, Union
import torch
@@ -1058,7 +1058,9 @@ class StableDiffusionControlNetPipeline(
torch.cuda.empty_cache()
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -1138,7 +1138,9 @@ class StableDiffusionControlNetImg2ImgPipeline(
torch.cuda.empty_cache()
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -1405,7 +1405,9 @@ class StableDiffusionControlNetInpaintPipeline(
torch.cuda.empty_cache()
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -1109,8 +1109,6 @@ class IFInpaintingSuperResolutionPipeline(DiffusionPipeline, LoraLoaderMixin):
nsfw_detected = None
watermark_detected = None
if hasattr(self, "unet_offload_hook") and self.unet_offload_hook is not None:
self.unet_offload_hook.offload()
else:
# 10. Post-processing
image = (image / 2 + 0.5).clamp(0, 1)
@@ -1119,9 +1117,7 @@ class IFInpaintingSuperResolutionPipeline(DiffusionPipeline, LoraLoaderMixin):
# 11. Run safety checker
image, nsfw_detected, watermark_detected = self.run_safety_checker(image, device, prompt_embeds.dtype)
# Offload last model to CPU
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.final_offload_hook.offload()
self.maybe_free_model_hooks()
if not return_dict:
return (image, nsfw_detected, watermark_detected)
@@ -388,6 +388,8 @@ class KandinskyPipeline(DiffusionPipeline):
# post-processing
image = self.movq.decode(latents, force_not_quantize=True)["sample"]
self.maybe_free_model_hooks()
if output_type not in ["pt", "np", "pil"]:
raise ValueError(f"Only the output types `pt`, `pil` and `np` are supported not output_type={output_type}")
@@ -321,6 +321,9 @@ class KandinskyCombinedPipeline(DiffusionPipeline):
callback_steps=callback_steps,
return_dict=return_dict,
)
self.maybe_free_model_hooks()
return outputs
@@ -558,6 +561,9 @@ class KandinskyImg2ImgCombinedPipeline(DiffusionPipeline):
callback_steps=callback_steps,
return_dict=return_dict,
)
self.maybe_free_model_hooks()
return outputs
@@ -593,7 +599,7 @@ class KandinskyInpaintCombinedPipeline(DiffusionPipeline):
"""
_load_connected_pipes = True
model_cpu_offload_seq = "prior_text_encoder->prior_image_encoder->prior_prior->" "text_encoder->unet->movq"
model_cpu_offload_seq = "prior_text_encoder->prior_image_encoder->prior_prior->text_encoder->unet->movq"
def __init__(
self,
@@ -802,4 +808,7 @@ class KandinskyInpaintCombinedPipeline(DiffusionPipeline):
callback_steps=callback_steps,
return_dict=return_dict,
)
self.maybe_free_model_hooks()
return outputs
@@ -481,6 +481,8 @@ class KandinskyImg2ImgPipeline(DiffusionPipeline):
# 7. post-processing
image = self.movq.decode(latents, force_not_quantize=True)["sample"]
self.maybe_free_model_hooks()
if output_type not in ["pt", "np", "pil"]:
raise ValueError(f"Only the output types `pt`, `pil` and `np` are supported not output_type={output_type}")
@@ -616,6 +616,8 @@ class KandinskyInpaintPipeline(DiffusionPipeline):
# post-processing
image = self.movq.decode(latents, force_not_quantize=True)["sample"]
self.maybe_free_model_hooks()
if output_type not in ["pt", "np", "pil"]:
raise ValueError(f"Only the output types `pt`, `pil` and `np` are supported not output_type={output_type}")
@@ -527,7 +527,7 @@ class KandinskyPriorPipeline(DiffusionPipeline):
if negative_prompt is None:
zero_embeds = self.get_zero_embed(latents.shape[0], device=latents.device)
self.maybe_free_model_hooks
self.maybe_free_model_hooks()
else:
image_embeddings, zero_embeds = image_embeddings.chunk(2)
@@ -326,6 +326,8 @@ class KandinskyV22CombinedPipeline(DiffusionPipeline):
callback_on_step_end=callback_on_step_end,
callback_on_step_end_tensor_inputs=callback_on_step_end_tensor_inputs,
)
self.maybe_free_model_hooks()
return outputs
@@ -572,6 +574,8 @@ class KandinskyV22Img2ImgCombinedPipeline(DiffusionPipeline):
callback_on_step_end=callback_on_step_end,
callback_on_step_end_tensor_inputs=callback_on_step_end_tensor_inputs,
)
self.maybe_free_model_hooks()
return outputs
@@ -842,4 +846,6 @@ class KandinskyV22InpaintCombinedPipeline(DiffusionPipeline):
callback_on_step_end_tensor_inputs=callback_on_step_end_tensor_inputs,
**kwargs,
)
self.maybe_free_model_hooks()
return outputs
@@ -531,14 +531,10 @@ class KandinskyV22PriorPipeline(DiffusionPipeline):
# if negative prompt has been defined, we retrieve split the image embedding into two
if negative_prompt is None:
zero_embeds = self.get_zero_embed(latents.shape[0], device=latents.device)
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.final_offload_hook.offload()
else:
image_embeddings, zero_embeds = image_embeddings.chunk(2)
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.prior_hook.offload()
self.maybe_free_model_hooks()
if output_type not in ["pt", "np"]:
raise ValueError(f"Only the output types `pt` and `np` are supported not output_type={output_type}")
@@ -545,12 +545,10 @@ class KandinskyV22PriorEmb2EmbPipeline(DiffusionPipeline):
# if negative prompt has been defined, we retrieve split the image embedding into two
if negative_prompt is None:
zero_embeds = self.get_zero_embed(latents.shape[0], device=latents.device)
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.final_offload_hook.offload()
else:
image_embeddings, zero_embeds = image_embeddings.chunk(2)
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.prior_hook.offload()
self.maybe_free_model_hooks()
if output_type not in ["pt", "np"]:
raise ValueError(f"Only the output types `pt` and `np` are supported not output_type={output_type}")
@@ -1,19 +1,40 @@
from typing import TYPE_CHECKING
from ...utils import (
DIFFUSERS_SLOW_IMPORT,
OptionalDependencyNotAvailable,
_LazyModule,
get_objects_from_module,
is_torch_available,
is_transformers_available,
)
_import_structure = {
"pipeline_latent_consistency_img2img": ["LatentConsistencyModelImg2ImgPipeline"],
"pipeline_latent_consistency_text2img": ["LatentConsistencyModelPipeline"],
}
_dummy_objects = {}
_import_structure = {}
if TYPE_CHECKING:
from .pipeline_latent_consistency_img2img import LatentConsistencyModelImg2ImgPipeline
from .pipeline_latent_consistency_text2img import LatentConsistencyModelPipeline
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_torch_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_torch_and_transformers_objects))
else:
_import_structure["pipeline_latent_consistency_img2img"] = ["LatentConsistencyModelImg2ImgPipeline"]
_import_structure["pipeline_latent_consistency_text2img"] = ["LatentConsistencyModelPipeline"]
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_torch_and_transformers_objects import *
else:
from .pipeline_latent_consistency_img2img import LatentConsistencyModelImg2ImgPipeline
from .pipeline_latent_consistency_text2img import LatentConsistencyModelPipeline
else:
import sys
@@ -24,3 +45,6 @@ else:
_import_structure,
module_spec=__spec__,
)
for name, value in _dummy_objects.items():
setattr(sys.modules[__name__], name, value)
@@ -60,7 +60,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> import PIL
>>> pipe = DiffusionPipeline.from_pretrained("SimianLuo/LCM_Dreamshaper_v7")
>>> pipe = AutoPipelineForImage2Image.from_pretrained("SimianLuo/LCM_Dreamshaper_v7")
>>> # To save GPU memory, torch.float16 can be used, but it may compromise image quality.
>>> pipe.to(torch_device="cuda", torch_dtype=torch.float32)
@@ -738,7 +738,7 @@ class LatentConsistencyModelImg2ImgPipeline(
if original_inference_steps is not None
else self.scheduler.config.original_inference_steps
)
latent_timestep = torch.tensor(int(strength * original_inference_steps))
latent_timestep = timesteps[:1]
latents = self.prepare_latents(
image, latent_timestep, batch_size, num_images_per_prompt, prompt_embeds.dtype, device, generator
)
+4 -5
View File
@@ -158,9 +158,9 @@ def is_safetensors_compatible(filenames, variant=None, passed_components=None) -
continue
if extension == ".bin":
pt_filenames.append(filename)
pt_filenames.append(os.path.normpath(filename))
elif extension == ".safetensors":
sf_filenames.add(filename)
sf_filenames.add(os.path.normpath(filename))
for filename in pt_filenames:
# filename = 'foo/bar/baz.bam' -> path = 'foo/bar', filename = 'baz', extention = '.bam'
@@ -172,9 +172,8 @@ def is_safetensors_compatible(filenames, variant=None, passed_components=None) -
else:
filename = filename
expected_sf_filename = os.path.join(path, filename)
expected_sf_filename = os.path.normpath(os.path.join(path, filename))
expected_sf_filename = f"{expected_sf_filename}.safetensors"
if expected_sf_filename not in sf_filenames:
logger.warning(f"{expected_sf_filename} not found")
return False
@@ -1774,7 +1773,7 @@ class DiffusionPipeline(ConfigMixin, PushToHubMixin):
)
):
raise EnvironmentError(
f"Could not found the necessary `safetensors` weights in {model_filenames} (variant={variant})"
f"Could not find the necessary `safetensors` weights in {model_filenames} (variant={variant})"
)
if from_flax:
ignore_patterns = ["*.bin", "*.safetensors", "*.onnx", "*.pb"]
@@ -1 +1,48 @@
from .pipeline_pixart_alpha import PixArtAlphaPipeline
from typing import TYPE_CHECKING
from ...utils import (
DIFFUSERS_SLOW_IMPORT,
OptionalDependencyNotAvailable,
_LazyModule,
get_objects_from_module,
is_torch_available,
is_transformers_available,
)
_dummy_objects = {}
_import_structure = {}
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_torch_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_torch_and_transformers_objects))
else:
_import_structure["pipeline_pixart_alpha"] = ["PixArtAlphaPipeline"]
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_torch_and_transformers_objects import *
else:
from .pipeline_pixart_alpha import PixArtAlphaPipeline
else:
import sys
sys.modules[__name__] = _LazyModule(
__name__,
globals()["__file__"],
_import_structure,
module_spec=__spec__,
)
for name, value in _dummy_objects.items():
setattr(sys.modules[__name__], name, value)
@@ -156,6 +156,8 @@ class PixArtAlphaPipeline(DiffusionPipeline):
mask_feature: (bool, defaults to `True`):
If `True`, the function will mask the text embeddings.
"""
embeds_initially_provided = prompt_embeds is not None and negative_prompt_embeds is not None
if device is None:
device = self._execution_device
@@ -253,7 +255,7 @@ class PixArtAlphaPipeline(DiffusionPipeline):
negative_prompt_embeds = None
# Perform additional masking.
if mask_feature:
if mask_feature and not embeds_initially_provided:
prompt_embeds = prompt_embeds.unsqueeze(1)
masked_prompt_embeds, keep_indices = self.mask_text_embeddings(prompt_embeds, prompt_embeds_attention_mask)
masked_prompt_embeds = masked_prompt_embeds.squeeze(1)
@@ -918,6 +918,7 @@ class CycleDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lor
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept]
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
@@ -576,6 +576,35 @@ class StableDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -593,7 +622,7 @@ class StableDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -790,6 +819,14 @@ class StableDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 6.5 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
# 7. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
@@ -804,6 +841,7 @@ class StableDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
return_dict=False,
)[0]
@@ -838,7 +876,9 @@ class StableDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Lo
callback(step_idx, t, latents)
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -1027,6 +1027,7 @@ class StableDiffusionAttendAndExcitePipeline(DiffusionPipeline, TextualInversion
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept]
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
@@ -846,6 +846,7 @@ class StableDiffusionDepth2ImgPipeline(DiffusionPipeline, TextualInversionLoader
image = latents
image = self.image_processor.postprocess(image, output_type=output_type)
self.maybe_free_model_hooks()
if not return_dict:
return (image,)
@@ -439,6 +439,8 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
@@ -640,6 +640,35 @@ class StableDiffusionImg2ImgPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -653,7 +682,7 @@ class StableDiffusionImg2ImgPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -841,6 +870,14 @@ class StableDiffusionImg2ImgPipeline(
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 7.5 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
# 8. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
@@ -855,6 +892,7 @@ class StableDiffusionImg2ImgPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
return_dict=False,
)[0]
@@ -885,7 +923,9 @@ class StableDiffusionImg2ImgPipeline(
callback(step_idx, t, latents)
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0]
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -765,6 +765,35 @@ class StableDiffusionInpaintPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -778,7 +807,7 @@ class StableDiffusionInpaintPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -1087,6 +1116,14 @@ class StableDiffusionInpaintPipeline(
# 9. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 9.5 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
# 10. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
@@ -1106,6 +1143,7 @@ class StableDiffusionInpaintPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
return_dict=False,
)[0]
@@ -1159,7 +1197,9 @@ class StableDiffusionInpaintPipeline(
init_image = self._encode_vae_image(init_image, generator=generator)
mask_condition = mask_condition.to(device=device, dtype=masked_image_latents.dtype)
condition_kwargs = {"image": init_image_condition, "mask": mask_condition}
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, **condition_kwargs)[0]
image = self.vae.decode(
latents / self.vae.config.scaling_factor, return_dict=False, generator=generator, **condition_kwargs
)[0]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
@@ -511,6 +511,8 @@ class StableDiffusionLatentUpscalePipeline(DiffusionPipeline, FromSingleFileMixi
image = self.image_processor.postprocess(image, output_type=output_type)
self.maybe_free_model_hooks()
if not return_dict:
return (image,)
@@ -802,6 +802,8 @@ class StableDiffusionPanoramaPipeline(DiffusionPipeline, TextualInversionLoaderM
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
@@ -741,6 +741,8 @@ class StableDiffusionSAGPipeline(DiffusionPipeline, TextualInversionLoaderMixin)
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
@@ -206,17 +206,15 @@ class StableUnCLIPPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
prior_text_encoder_output = self.prior_text_encoder(text_input_ids.to(device))
prompt_embeds = prior_text_encoder_output.text_embeds
prior_text_encoder_hidden_states = prior_text_encoder_output.last_hidden_state
text_enc_hid_states = prior_text_encoder_output.last_hidden_state
else:
batch_size = text_model_output[0].shape[0]
prompt_embeds, prior_text_encoder_hidden_states = text_model_output[0], text_model_output[1]
prompt_embeds, text_enc_hid_states = text_model_output[0], text_model_output[1]
text_mask = text_attention_mask
prompt_embeds = prompt_embeds.repeat_interleave(num_images_per_prompt, dim=0)
prior_text_encoder_hidden_states = prior_text_encoder_hidden_states.repeat_interleave(
num_images_per_prompt, dim=0
)
text_enc_hid_states = text_enc_hid_states.repeat_interleave(num_images_per_prompt, dim=0)
text_mask = text_mask.repeat_interleave(num_images_per_prompt, dim=0)
if do_classifier_free_guidance:
@@ -235,9 +233,7 @@ class StableUnCLIPPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
)
negative_prompt_embeds = negative_prompt_embeds_prior_text_encoder_output.text_embeds
uncond_prior_text_encoder_hidden_states = (
negative_prompt_embeds_prior_text_encoder_output.last_hidden_state
)
uncond_text_enc_hid_states = negative_prompt_embeds_prior_text_encoder_output.last_hidden_state
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
@@ -245,11 +241,9 @@ class StableUnCLIPPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len)
seq_len = uncond_prior_text_encoder_hidden_states.shape[1]
uncond_prior_text_encoder_hidden_states = uncond_prior_text_encoder_hidden_states.repeat(
1, num_images_per_prompt, 1
)
uncond_prior_text_encoder_hidden_states = uncond_prior_text_encoder_hidden_states.view(
seq_len = uncond_text_enc_hid_states.shape[1]
uncond_text_enc_hid_states = uncond_text_enc_hid_states.repeat(1, num_images_per_prompt, 1)
uncond_text_enc_hid_states = uncond_text_enc_hid_states.view(
batch_size * num_images_per_prompt, seq_len, -1
)
uncond_text_mask = uncond_text_mask.repeat_interleave(num_images_per_prompt, dim=0)
@@ -260,13 +254,11 @@ class StableUnCLIPPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraL
# Here we concatenate the unconditional and text embeddings into a single batch
# to avoid doing two forward passes
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds])
prior_text_encoder_hidden_states = torch.cat(
[uncond_prior_text_encoder_hidden_states, prior_text_encoder_hidden_states]
)
text_enc_hid_states = torch.cat([uncond_text_enc_hid_states, text_enc_hid_states])
text_mask = torch.cat([uncond_text_mask, text_mask])
return prompt_embeds, prior_text_encoder_hidden_states, text_mask
return prompt_embeds, text_enc_hid_states, text_mask
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._encode_prompt
def _encode_prompt(
@@ -636,6 +636,35 @@ class StableDiffusionXLPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -653,7 +682,7 @@ class StableDiffusionXLPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -989,6 +1018,14 @@ class StableDiffusionXLPipeline(
num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps)))
timesteps = timesteps[:num_inference_steps]
# 9. Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
self._num_timesteps = len(timesteps)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
@@ -1003,6 +1040,7 @@ class StableDiffusionXLPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
added_cond_kwargs=added_cond_kwargs,
return_dict=False,
@@ -763,6 +763,35 @@ class StableDiffusionXLImg2ImgPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -780,7 +809,7 @@ class StableDiffusionXLImg2ImgPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -1156,6 +1185,15 @@ class StableDiffusionXLImg2ImgPipeline(
)
num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps)))
timesteps = timesteps[:num_inference_steps]
# 9.2 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
self._num_timesteps = len(timesteps)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
@@ -1170,6 +1208,7 @@ class StableDiffusionXLImg2ImgPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
added_cond_kwargs=added_cond_kwargs,
return_dict=False,
@@ -982,6 +982,35 @@ class StableDiffusionXLInpaintPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
@property
def guidance_scale(self):
return self._guidance_scale
@@ -999,7 +1028,7 @@ class StableDiffusionXLInpaintPipeline(
# corresponds to doing no classifier free guidance.
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
@property
def cross_attention_kwargs(self):
@@ -1464,6 +1493,14 @@ class StableDiffusionXLInpaintPipeline(
num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps)))
timesteps = timesteps[:num_inference_steps]
# 11.1 Optionally get Guidance Scale Embedding
timestep_cond = None
if self.unet.config.time_cond_proj_dim is not None:
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt)
timestep_cond = self.get_guidance_scale_embedding(
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim
).to(device=device, dtype=latents.dtype)
self._num_timesteps = len(timesteps)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
@@ -1482,6 +1519,7 @@ class StableDiffusionXLInpaintPipeline(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
timestep_cond=timestep_cond,
cross_attention_kwargs=self.cross_attention_kwargs,
added_cond_kwargs=added_cond_kwargs,
return_dict=False,
@@ -156,15 +156,15 @@ class UnCLIPPipeline(DiffusionPipeline):
text_encoder_output = self.text_encoder(text_input_ids.to(device))
prompt_embeds = text_encoder_output.text_embeds
text_encoder_hidden_states = text_encoder_output.last_hidden_state
text_enc_hid_states = text_encoder_output.last_hidden_state
else:
batch_size = text_model_output[0].shape[0]
prompt_embeds, text_encoder_hidden_states = text_model_output[0], text_model_output[1]
prompt_embeds, text_enc_hid_states = text_model_output[0], text_model_output[1]
text_mask = text_attention_mask
prompt_embeds = prompt_embeds.repeat_interleave(num_images_per_prompt, dim=0)
text_encoder_hidden_states = text_encoder_hidden_states.repeat_interleave(num_images_per_prompt, dim=0)
text_enc_hid_states = text_enc_hid_states.repeat_interleave(num_images_per_prompt, dim=0)
text_mask = text_mask.repeat_interleave(num_images_per_prompt, dim=0)
if do_classifier_free_guidance:
@@ -181,7 +181,7 @@ class UnCLIPPipeline(DiffusionPipeline):
negative_prompt_embeds_text_encoder_output = self.text_encoder(uncond_input.input_ids.to(device))
negative_prompt_embeds = negative_prompt_embeds_text_encoder_output.text_embeds
uncond_text_encoder_hidden_states = negative_prompt_embeds_text_encoder_output.last_hidden_state
uncond_text_enc_hid_states = negative_prompt_embeds_text_encoder_output.last_hidden_state
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
@@ -189,9 +189,9 @@ class UnCLIPPipeline(DiffusionPipeline):
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len)
seq_len = uncond_text_encoder_hidden_states.shape[1]
uncond_text_encoder_hidden_states = uncond_text_encoder_hidden_states.repeat(1, num_images_per_prompt, 1)
uncond_text_encoder_hidden_states = uncond_text_encoder_hidden_states.view(
seq_len = uncond_text_enc_hid_states.shape[1]
uncond_text_enc_hid_states = uncond_text_enc_hid_states.repeat(1, num_images_per_prompt, 1)
uncond_text_enc_hid_states = uncond_text_enc_hid_states.view(
batch_size * num_images_per_prompt, seq_len, -1
)
uncond_text_mask = uncond_text_mask.repeat_interleave(num_images_per_prompt, dim=0)
@@ -202,11 +202,11 @@ class UnCLIPPipeline(DiffusionPipeline):
# Here we concatenate the unconditional and text embeddings into a single batch
# to avoid doing two forward passes
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds])
text_encoder_hidden_states = torch.cat([uncond_text_encoder_hidden_states, text_encoder_hidden_states])
text_enc_hid_states = torch.cat([uncond_text_enc_hid_states, text_enc_hid_states])
text_mask = torch.cat([uncond_text_mask, text_mask])
return prompt_embeds, text_encoder_hidden_states, text_mask
return prompt_embeds, text_enc_hid_states, text_mask
@torch.no_grad()
def __call__(
@@ -293,7 +293,7 @@ class UnCLIPPipeline(DiffusionPipeline):
do_classifier_free_guidance = prior_guidance_scale > 1.0 or decoder_guidance_scale > 1.0
prompt_embeds, text_encoder_hidden_states, text_mask = self._encode_prompt(
prompt_embeds, text_enc_hid_states, text_mask = self._encode_prompt(
prompt, device, num_images_per_prompt, do_classifier_free_guidance, text_model_output, text_attention_mask
)
@@ -321,7 +321,7 @@ class UnCLIPPipeline(DiffusionPipeline):
latent_model_input,
timestep=t,
proj_embedding=prompt_embeds,
encoder_hidden_states=text_encoder_hidden_states,
encoder_hidden_states=text_enc_hid_states,
attention_mask=text_mask,
).predicted_image_embedding
@@ -352,10 +352,10 @@ class UnCLIPPipeline(DiffusionPipeline):
# decoder
text_encoder_hidden_states, additive_clip_time_embeddings = self.text_proj(
text_enc_hid_states, additive_clip_time_embeddings = self.text_proj(
image_embeddings=image_embeddings,
prompt_embeds=prompt_embeds,
text_encoder_hidden_states=text_encoder_hidden_states,
text_encoder_hidden_states=text_enc_hid_states,
do_classifier_free_guidance=do_classifier_free_guidance,
)
@@ -377,7 +377,7 @@ class UnCLIPPipeline(DiffusionPipeline):
decoder_latents = self.prepare_latents(
(batch_size, num_channels_latents, height, width),
text_encoder_hidden_states.dtype,
text_enc_hid_states.dtype,
device,
generator,
decoder_latents,
@@ -391,7 +391,7 @@ class UnCLIPPipeline(DiffusionPipeline):
noise_pred = self.decoder(
sample=latent_model_input,
timestep=t,
encoder_hidden_states=text_encoder_hidden_states,
encoder_hidden_states=text_enc_hid_states,
class_labels=additive_clip_time_embeddings,
attention_mask=decoder_text_mask,
).sample
@@ -1494,7 +1494,6 @@ class ResnetBlockFlat(nn.Module):
return output_tensor
# Copied from diffusers.models.unet_2d_blocks.DownBlock2D with DownBlock2D->DownBlockFlat, ResnetBlock2D->ResnetBlockFlat, Downsample2D->LinearMultiDim
class DownBlockFlat(nn.Module):
def __init__(
self,
@@ -1583,7 +1582,6 @@ class DownBlockFlat(nn.Module):
return hidden_states, output_states
# Copied from diffusers.models.unet_2d_blocks.CrossAttnDownBlock2D with CrossAttnDownBlock2D->CrossAttnDownBlockFlat, ResnetBlock2D->ResnetBlockFlat, Downsample2D->LinearMultiDim
class CrossAttnDownBlockFlat(nn.Module):
def __init__(
self,
+2
View File
@@ -38,6 +38,7 @@ except OptionalDependencyNotAvailable:
_dummy_modules.update(get_objects_from_module(dummy_pt_objects))
else:
_import_structure["scheduling_consistency_decoder"] = ["ConsistencyDecoderScheduler"]
_import_structure["scheduling_consistency_models"] = ["CMStochasticIterativeScheduler"]
_import_structure["scheduling_ddim"] = ["DDIMScheduler"]
_import_structure["scheduling_ddim_inverse"] = ["DDIMInverseScheduler"]
@@ -128,6 +129,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
except OptionalDependencyNotAvailable:
from ..utils.dummy_pt_objects import * # noqa F403
else:
from .scheduling_consistency_decoder import ConsistencyDecoderScheduler
from .scheduling_consistency_models import CMStochasticIterativeScheduler
from .scheduling_ddim import DDIMScheduler
from .scheduling_ddim_inverse import DDIMInverseScheduler
@@ -0,0 +1,180 @@
import math
from dataclasses import dataclass
from typing import Optional, Tuple, Union
import torch
from ..configuration_utils import ConfigMixin, register_to_config
from ..utils import BaseOutput
from ..utils.torch_utils import randn_tensor
from .scheduling_utils import SchedulerMixin
# Copied from diffusers.schedulers.scheduling_ddpm.betas_for_alpha_bar
def betas_for_alpha_bar(
num_diffusion_timesteps,
max_beta=0.999,
alpha_transform_type="cosine",
):
"""
Create a beta schedule that discretizes the given alpha_t_bar function, which defines the cumulative product of
(1-beta) over time from t = [0,1].
Contains a function alpha_bar that takes an argument t and transforms it to the cumulative product of (1-beta) up
to that part of the diffusion process.
Args:
num_diffusion_timesteps (`int`): the number of betas to produce.
max_beta (`float`): the maximum beta to use; use values lower than 1 to
prevent singularities.
alpha_transform_type (`str`, *optional*, default to `cosine`): the type of noise schedule for alpha_bar.
Choose from `cosine` or `exp`
Returns:
betas (`np.ndarray`): the betas used by the scheduler to step the model outputs
"""
if alpha_transform_type == "cosine":
def alpha_bar_fn(t):
return math.cos((t + 0.008) / 1.008 * math.pi / 2) ** 2
elif alpha_transform_type == "exp":
def alpha_bar_fn(t):
return math.exp(t * -12.0)
else:
raise ValueError(f"Unsupported alpha_tranform_type: {alpha_transform_type}")
betas = []
for i in range(num_diffusion_timesteps):
t1 = i / num_diffusion_timesteps
t2 = (i + 1) / num_diffusion_timesteps
betas.append(min(1 - alpha_bar_fn(t2) / alpha_bar_fn(t1), max_beta))
return torch.tensor(betas, dtype=torch.float32)
@dataclass
class ConsistencyDecoderSchedulerOutput(BaseOutput):
"""
Output class for the scheduler's `step` function.
Args:
prev_sample (`torch.FloatTensor` of shape `(batch_size, num_channels, height, width)` for images):
Computed sample `(x_{t-1})` of previous timestep. `prev_sample` should be used as next model input in the
denoising loop.
"""
prev_sample: torch.FloatTensor
class ConsistencyDecoderScheduler(SchedulerMixin, ConfigMixin):
order = 1
@register_to_config
def __init__(
self,
num_train_timesteps: int = 1024,
sigma_data: float = 0.5,
):
betas = betas_for_alpha_bar(num_train_timesteps)
alphas = 1.0 - betas
alphas_cumprod = torch.cumprod(alphas, dim=0)
self.sqrt_alphas_cumprod = torch.sqrt(alphas_cumprod)
self.sqrt_one_minus_alphas_cumprod = torch.sqrt(1.0 - alphas_cumprod)
sigmas = torch.sqrt(1.0 / alphas_cumprod - 1)
sqrt_recip_alphas_cumprod = torch.sqrt(1.0 / alphas_cumprod)
self.c_skip = sqrt_recip_alphas_cumprod * sigma_data**2 / (sigmas**2 + sigma_data**2)
self.c_out = sigmas * sigma_data / (sigmas**2 + sigma_data**2) ** 0.5
self.c_in = sqrt_recip_alphas_cumprod / (sigmas**2 + sigma_data**2) ** 0.5
def set_timesteps(
self,
num_inference_steps: Optional[int] = None,
device: Union[str, torch.device] = None,
):
if num_inference_steps != 2:
raise ValueError("Currently more than 2 inference steps are not supported.")
self.timesteps = torch.tensor([1008, 512], dtype=torch.long, device=device)
self.sqrt_alphas_cumprod = self.sqrt_alphas_cumprod.to(device)
self.sqrt_one_minus_alphas_cumprod = self.sqrt_one_minus_alphas_cumprod.to(device)
self.c_skip = self.c_skip.to(device)
self.c_out = self.c_out.to(device)
self.c_in = self.c_in.to(device)
@property
def init_noise_sigma(self):
return self.sqrt_one_minus_alphas_cumprod[self.timesteps[0]]
def scale_model_input(self, sample: torch.FloatTensor, timestep: Optional[int] = None) -> torch.FloatTensor:
"""
Ensures interchangeability with schedulers that need to scale the denoising model input depending on the
current timestep.
Args:
sample (`torch.FloatTensor`):
The input sample.
timestep (`int`, *optional*):
The current timestep in the diffusion chain.
Returns:
`torch.FloatTensor`:
A scaled input sample.
"""
return sample * self.c_in[timestep]
def step(
self,
model_output: torch.FloatTensor,
timestep: Union[float, torch.FloatTensor],
sample: torch.FloatTensor,
generator: Optional[torch.Generator] = None,
return_dict: bool = True,
) -> Union[ConsistencyDecoderSchedulerOutput, Tuple]:
"""
Predict the sample from the previous timestep by reversing the SDE. This function propagates the diffusion
process from the learned model outputs (most often the predicted noise).
Args:
model_output (`torch.FloatTensor`):
The direct output from the learned diffusion model.
timestep (`float`):
The current timestep in the diffusion chain.
sample (`torch.FloatTensor`):
A current instance of a sample created by the diffusion process.
generator (`torch.Generator`, *optional*):
A random number generator.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a
[`~schedulers.scheduling_consistency_models.ConsistencyDecoderSchedulerOutput`] or `tuple`.
Returns:
[`~schedulers.scheduling_consistency_models.ConsistencyDecoderSchedulerOutput`] or `tuple`:
If return_dict is `True`,
[`~schedulers.scheduling_consistency_models.ConsistencyDecoderSchedulerOutput`] is returned, otherwise
a tuple is returned where the first element is the sample tensor.
"""
x_0 = self.c_out[timestep] * model_output + self.c_skip[timestep] * sample
timestep_idx = torch.where(self.timesteps == timestep)[0]
if timestep_idx == len(self.timesteps) - 1:
prev_sample = x_0
else:
noise = randn_tensor(x_0.shape, generator=generator, dtype=x_0.dtype, device=x_0.device)
prev_sample = (
self.sqrt_alphas_cumprod[self.timesteps[timestep_idx + 1]].to(x_0.dtype) * x_0
+ self.sqrt_one_minus_alphas_cumprod[self.timesteps[timestep_idx + 1]].to(x_0.dtype) * noise
)
if not return_dict:
return (prev_sample,)
return ConsistencyDecoderSchedulerOutput(prev_sample=prev_sample)
+15 -7
View File
@@ -182,6 +182,10 @@ class LCMScheduler(SchedulerMixin, ConfigMixin):
timestep_spacing (`str`, defaults to `"leading"`):
The way the timesteps should be scaled. Refer to Table 2 of the [Common Diffusion Noise Schedules and
Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) for more information.
timestep_scaling (`float`, defaults to 10.0):
The factor the timesteps will be multiplied by when calculating the consistency model boundary conditions
`c_skip` and `c_out`. Increasing this will decrease the approximation error (although the approximation
error at the default of `10.0` is already pretty small).
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
@@ -208,6 +212,7 @@ class LCMScheduler(SchedulerMixin, ConfigMixin):
dynamic_thresholding_ratio: float = 0.995,
sample_max_value: float = 1.0,
timestep_spacing: str = "leading",
timestep_scaling: float = 10.0,
rescale_betas_zero_snr: bool = False,
):
if trained_betas is not None:
@@ -380,12 +385,12 @@ class LCMScheduler(SchedulerMixin, ConfigMixin):
self._step_index = None
def get_scalings_for_boundary_condition_discrete(self, t):
def get_scalings_for_boundary_condition_discrete(self, timestep):
self.sigma_data = 0.5 # Default: 0.5
scaled_timestep = timestep * self.config.timestep_scaling
# By dividing 0.1: This is almost a delta function at t=0.
c_skip = self.sigma_data**2 / ((t / 0.1) ** 2 + self.sigma_data**2)
c_out = (t / 0.1) / ((t / 0.1) ** 2 + self.sigma_data**2) ** 0.5
c_skip = self.sigma_data**2 / (scaled_timestep**2 + self.sigma_data**2)
c_out = scaled_timestep / (scaled_timestep**2 + self.sigma_data**2) ** 0.5
return c_skip, c_out
def step(
@@ -466,9 +471,12 @@ class LCMScheduler(SchedulerMixin, ConfigMixin):
denoised = c_out * predicted_original_sample + c_skip * sample
# 7. Sample and inject noise z ~ N(0, I) for MultiStep Inference
# Noise is not used for one-step sampling.
if len(self.timesteps) > 1:
noise = randn_tensor(model_output.shape, generator=generator, device=model_output.device)
# Noise is not used on the final timestep of the timestep schedule.
# This also means that noise is not used for one-step sampling.
if self.step_index != self.num_inference_steps - 1:
noise = randn_tensor(
model_output.shape, generator=generator, device=model_output.device, dtype=denoised.dtype
)
prev_sample = alpha_prod_t_prev.sqrt() * denoised + beta_prod_t_prev.sqrt() * noise
else:
prev_sample = denoised
+4
View File
@@ -18,6 +18,7 @@ from huggingface_hub.constants import HUGGINGFACE_HUB_CACHE, hf_cache_home
from packaging import version
from .import_utils import is_peft_available, is_transformers_available
from ..dependency_versions_check import dep_version_check
default_cache_path = HUGGINGFACE_HUB_CACHE
@@ -50,3 +51,6 @@ _required_transformers_version = is_transformers_available() and version.parse(
) > version.parse(MIN_TRANSFORMERS_VERSION)
USE_PEFT_BACKEND = _required_peft_version and _required_transformers_version
if USE_PEFT_BACKEND:
dep_version_check("peft")
+15
View File
@@ -47,6 +47,21 @@ class AutoencoderTiny(metaclass=DummyObject):
requires_backends(cls, ["torch"])
class ConsistencyDecoderVAE(metaclass=DummyObject):
_backends = ["torch"]
def __init__(self, *args, **kwargs):
requires_backends(self, ["torch"])
@classmethod
def from_config(cls, *args, **kwargs):
requires_backends(cls, ["torch"])
@classmethod
def from_pretrained(cls, *args, **kwargs):
requires_backends(cls, ["torch"])
class ControlNetModel(metaclass=DummyObject):
_backends = ["torch"]
+117
View File
@@ -0,0 +1,117 @@
# Copyright 2020 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
"""
Utilities for working with package versions
"""
import importlib.metadata
import operator
import re
import sys
from typing import Optional
from packaging import version
ops = {
"<": operator.lt,
"<=": operator.le,
"==": operator.eq,
"!=": operator.ne,
">=": operator.ge,
">": operator.gt,
}
def _compare_versions(op, got_ver, want_ver, requirement, pkg, hint):
if got_ver is None or want_ver is None:
raise ValueError(
f"Unable to compare versions for {requirement}: need={want_ver} found={got_ver}. This is unusual. Consider"
f" reinstalling {pkg}."
)
if not ops[op](version.parse(got_ver), version.parse(want_ver)):
raise ImportError(
f"{requirement} is required for a normal functioning of this module, but found {pkg}=={got_ver}.{hint}"
)
def require_version(requirement: str, hint: Optional[str] = None) -> None:
"""
Perform a runtime check of the dependency versions, using the exact same syntax used by pip.
The installed module version comes from the *site-packages* dir via *importlib.metadata*.
Args:
requirement (`str`): pip style definition, e.g., "tokenizers==0.9.4", "tqdm>=4.27", "numpy"
hint (`str`, *optional*): what suggestion to print in case of requirements not being met
Example:
```python
require_version("pandas>1.1.2")
require_version("numpy>1.18.5", "this is important to have for whatever reason")
```"""
hint = f"\n{hint}" if hint is not None else ""
# non-versioned check
if re.match(r"^[\w_\-\d]+$", requirement):
pkg, op, want_ver = requirement, None, None
else:
match = re.findall(r"^([^!=<>\s]+)([\s!=<>]{1,2}.+)", requirement)
if not match:
raise ValueError(
"requirement needs to be in the pip package format, .e.g., package_a==1.23, or package_b>=1.23, but"
f" got {requirement}"
)
pkg, want_full = match[0]
want_range = want_full.split(",") # there could be multiple requirements
wanted = {}
for w in want_range:
match = re.findall(r"^([\s!=<>]{1,2})(.+)", w)
if not match:
raise ValueError(
"requirement needs to be in the pip package format, .e.g., package_a==1.23, or package_b>=1.23,"
f" but got {requirement}"
)
op, want_ver = match[0]
wanted[op] = want_ver
if op not in ops:
raise ValueError(f"{requirement}: need one of {list(ops.keys())}, but got {op}")
# special case
if pkg == "python":
got_ver = ".".join([str(x) for x in sys.version_info[:3]])
for op, want_ver in wanted.items():
_compare_versions(op, got_ver, want_ver, requirement, pkg, hint)
return
# check if any version is installed
try:
got_ver = importlib.metadata.version(pkg)
except importlib.metadata.PackageNotFoundError:
raise importlib.metadata.PackageNotFoundError(
f"The '{requirement}' distribution was not found and is required by this application. {hint}"
)
# check that the right version is installed if version number or a range was provided
if want_ver is not None:
for op, want_ver in wanted.items():
_compare_versions(op, got_ver, want_ver, requirement, pkg, hint)
def require_version_core(requirement):
"""require_version wrapper which emits a core-specific hint on failure"""
hint = "Try: pip install transformers -U or pip install -e '.[dev]' if you're working with git main"
return require_version(requirement, hint)

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