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@@ -202,7 +202,6 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
|
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|
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- https://github.com/microsoft/TaskMatrix
|
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- https://github.com/invoke-ai/InvokeAI
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- https://github.com/InstantID/InstantID
|
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- https://github.com/apple/ml-stable-diffusion
|
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- https://github.com/Sanster/lama-cleaner
|
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- https://github.com/IDEA-Research/Grounded-Segment-Anything
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|
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+62
-70
@@ -223,76 +223,68 @@
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sections:
|
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- local: api/models/overview
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title: Overview
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- sections:
|
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- local: api/models/controlnet
|
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title: ControlNetModel
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- local: api/models/controlnet_hunyuandit
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title: HunyuanDiT2DControlNetModel
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- local: api/models/controlnet_sd3
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title: SD3ControlNetModel
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- local: api/models/controlnet_sparsectrl
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title: SparseControlNetModel
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title: ControlNets
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- sections:
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- local: api/models/aura_flow_transformer2d
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title: AuraFlowTransformer2DModel
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- local: api/models/cogvideox_transformer3d
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title: CogVideoXTransformer3DModel
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- local: api/models/dit_transformer2d
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title: DiTTransformer2DModel
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- local: api/models/flux_transformer
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title: FluxTransformer2DModel
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- local: api/models/hunyuan_transformer2d
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title: HunyuanDiT2DModel
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- local: api/models/latte_transformer3d
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title: LatteTransformer3DModel
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- local: api/models/lumina_nextdit2d
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title: LuminaNextDiT2DModel
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- local: api/models/pixart_transformer2d
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title: PixArtTransformer2DModel
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- local: api/models/prior_transformer
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title: PriorTransformer
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- local: api/models/sd3_transformer2d
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title: SD3Transformer2DModel
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- local: api/models/stable_audio_transformer
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title: StableAudioDiTModel
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- local: api/models/transformer2d
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title: Transformer2DModel
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- local: api/models/transformer_temporal
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title: TransformerTemporalModel
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title: Transformers
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- sections:
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- local: api/models/stable_cascade_unet
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title: StableCascadeUNet
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- local: api/models/unet
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title: UNet1DModel
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- local: api/models/unet2d
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title: UNet2DModel
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- local: api/models/unet2d-cond
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title: UNet2DConditionModel
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- local: api/models/unet3d-cond
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title: UNet3DConditionModel
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- local: api/models/unet-motion
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title: UNetMotionModel
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- local: api/models/uvit2d
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title: UViT2DModel
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title: UNets
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- sections:
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- local: api/models/autoencoderkl
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title: AutoencoderKL
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- local: api/models/autoencoderkl_cogvideox
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title: AutoencoderKLCogVideoX
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- local: api/models/asymmetricautoencoderkl
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title: AsymmetricAutoencoderKL
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- local: api/models/consistency_decoder_vae
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title: ConsistencyDecoderVAE
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- local: api/models/autoencoder_oobleck
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title: Oobleck AutoEncoder
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- local: api/models/autoencoder_tiny
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title: Tiny AutoEncoder
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- local: api/models/vq
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title: VQModel
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title: VAEs
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- local: api/models/unet
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title: UNet1DModel
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- local: api/models/unet2d
|
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title: UNet2DModel
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- local: api/models/unet2d-cond
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title: UNet2DConditionModel
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- local: api/models/unet3d-cond
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title: UNet3DConditionModel
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- local: api/models/unet-motion
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title: UNetMotionModel
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- local: api/models/uvit2d
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title: UViT2DModel
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- local: api/models/vq
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title: VQModel
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- local: api/models/autoencoderkl
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title: AutoencoderKL
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- local: api/models/autoencoderkl_cogvideox
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title: AutoencoderKLCogVideoX
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- local: api/models/asymmetricautoencoderkl
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title: AsymmetricAutoencoderKL
|
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- local: api/models/stable_cascade_unet
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title: StableCascadeUNet
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- local: api/models/autoencoder_tiny
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title: Tiny AutoEncoder
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- local: api/models/autoencoder_oobleck
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title: Oobleck AutoEncoder
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- local: api/models/consistency_decoder_vae
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title: ConsistencyDecoderVAE
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- local: api/models/transformer2d
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title: Transformer2DModel
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- local: api/models/pixart_transformer2d
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title: PixArtTransformer2DModel
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- local: api/models/dit_transformer2d
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title: DiTTransformer2DModel
|
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- local: api/models/hunyuan_transformer2d
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title: HunyuanDiT2DModel
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- local: api/models/aura_flow_transformer2d
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title: AuraFlowTransformer2DModel
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- local: api/models/flux_transformer
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title: FluxTransformer2DModel
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- local: api/models/latte_transformer3d
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title: LatteTransformer3DModel
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- local: api/models/cogvideox_transformer3d
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title: CogVideoXTransformer3DModel
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- local: api/models/lumina_nextdit2d
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title: LuminaNextDiT2DModel
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- local: api/models/transformer_temporal
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title: TransformerTemporalModel
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- local: api/models/sd3_transformer2d
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title: SD3Transformer2DModel
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- local: api/models/stable_audio_transformer
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title: StableAudioDiTModel
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- local: api/models/prior_transformer
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title: PriorTransformer
|
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- local: api/models/controlnet
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title: ControlNetModel
|
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- local: api/models/controlnet_hunyuandit
|
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title: HunyuanDiT2DControlNetModel
|
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- local: api/models/controlnet_sd3
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title: SD3ControlNetModel
|
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- local: api/models/controlnet_sparsectrl
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title: SparseControlNetModel
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title: Models
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- isExpanded: false
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sections:
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@@ -15,9 +15,7 @@
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# CogVideoX
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<!-- TODO: update paper with ArXiv link when ready. -->
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[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://github.com/THUDM/CogVideo/blob/main/resources/CogVideoX.pdf) from Tsinghua University & ZhipuAI.
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[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
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|
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The abstract from the paper is:
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|
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@@ -31,6 +29,10 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.m
|
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|
||||
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
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There are two models available that can be used with the CogVideoX pipeline:
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- [`THUDM/CogVideoX-2b`](https://huggingface.co/THUDM/CogVideoX-2b)
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- [`THUDM/CogVideoX-5b`](https://huggingface.co/THUDM/CogVideoX-5b)
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## Inference
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Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) to reduce the inference latency.
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@@ -43,43 +45,42 @@ from diffusers import CogVideoXPipeline
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from diffusers.utils import export_to_video
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pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-2b").to("cuda")
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prompt = (
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"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
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"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
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"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
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"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
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"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
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"atmosphere of this unique musical performance."
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)
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video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
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export_to_video(video, "output.mp4", fps=8)
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```
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Then change the memory layout of the pipelines `transformer` and `vae` components to `torch.channels-last`:
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Then change the memory layout of the pipelines `transformer` component to `torch.channels_last`:
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||||
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||||
```python
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pipeline.transformer.to(memory_format=torch.channels_last)
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pipeline.vae.to(memory_format=torch.channels_last)
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pipe.transformer.to(memory_format=torch.channels_last)
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```
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Finally, compile the components and run inference:
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```python
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pipeline.transformer = torch.compile(pipeline.transformer)
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pipeline.vae.decode = torch.compile(pipeline.vae.decode)
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pipe.transformer = torch.compile(pipeline.transformer, mode="max-autotune", fullgraph=True)
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|
||||
# CogVideoX works very well with long and well-described prompts
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# CogVideoX works well with long and well-described prompts
|
||||
prompt = "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical atmosphere of this unique musical performance."
|
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video = pipeline(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
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video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
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```
|
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||||
The [benchmark](TODO: link) results on an 80GB A100 machine are:
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The [benchmark](https://gist.github.com/a-r-r-o-w/5183d75e452a368fd17448fcc810bd3f) results on an 80GB A100 machine are:
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|
||||
```
|
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Without torch.compile(): Average inference time: TODO seconds.
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With torch.compile(): Average inference time: TODO seconds.
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Without torch.compile(): Average inference time: 96.89 seconds.
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With torch.compile(): Average inference time: 76.27 seconds.
|
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```
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### Memory optimization
|
||||
|
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CogVideoX-2b requires about 19 GB of GPU memory to decode 49 frames (6 seconds of video at 8 FPS) with output resolution 720x480 (W x H), which makes it not possible to run on consumer GPUs or free-tier T4 Colab. The following memory optimizations could be used to reduce the memory footprint. For replication, you can refer to [this](https://gist.github.com/a-r-r-o-w/3959a03f15be5c9bd1fe545b09dfcc93) script.
|
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|
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- `pipe.enable_model_cpu_offload()`:
|
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- Without enabling cpu offloading, memory usage is `33 GB`
|
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- With enabling cpu offloading, memory usage is `19 GB`
|
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- `pipe.vae.enable_tiling()`:
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- With enabling cpu offloading and tiling, memory usage is `11 GB`
|
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- `pipe.vae.enable_slicing()`
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|
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## CogVideoXPipeline
|
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|
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[[autodoc]] CogVideoXPipeline
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|
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@@ -21,7 +21,7 @@ Stable Audio is trained on a corpus of around 48k audio recordings, where around
|
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The abstract of the paper is the following:
|
||||
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
|
||||
|
||||
This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tools](https://github.com/Stability-AI/stable-audio-tools).
|
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This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tool](https://github.com/Stability-AI/stable-audio-tool).
|
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|
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## Tips
|
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|
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|
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@@ -125,5 +125,3 @@ image
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<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion + Tiny AutoEncoder</figcaption>
|
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</div>
|
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</div>
|
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|
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More tiny autoencoder models for other Stable Diffusion models, like Stable Diffusion 3, are available from [madebyollin](https://huggingface.co/madebyollin).
|
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@@ -14,9 +14,9 @@ specific language governing permissions and limitations under the License.
|
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|
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It can be fun and creative to use multiple [LoRAs]((https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora)) together to generate something entirely new and unique. This works by merging multiple LoRA weights together to produce images that are a blend of different styles. Diffusers provides a few methods to merge LoRAs depending on *how* you want to merge their weights, which can affect image quality.
|
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|
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This guide will show you how to merge LoRAs using the [`~loaders.PeftAdapterMixin.set_adapters`] and [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
|
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This guide will show you how to merge LoRAs using the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
|
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|
||||
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style](https://huggingface.co/KappaNeuro/studio-ghibli-style) and [Norod78/sdxl-chalkboarddrawing-lora](https://huggingface.co/Norod78/sdxl-chalkboarddrawing-lora) LoRAs with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
|
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For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style]() and [Norod78/sdxl-chalkboarddrawing-lora]() LoRAs with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
|
||||
|
||||
```py
|
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from diffusers import DiffusionPipeline
|
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@@ -29,7 +29,7 @@ pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_
|
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|
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## set_adapters
|
||||
|
||||
The [`~loaders.PeftAdapterMixin.set_adapters`] method merges LoRA adapters by concatenating their weighted matrices. Use the adapter name to specify which LoRAs to merge, and the `adapter_weights` parameter to control the scaling for each LoRA. For example, if `adapter_weights=[0.5, 0.5]`, then the merged LoRA output is an average of both LoRAs. Try adjusting the adapter weights to see how it affects the generated image!
|
||||
The [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method merges LoRA adapters by concatenating their weighted matrices. Use the adapter name to specify which LoRAs to merge, and the `adapter_weights` parameter to control the scaling for each LoRA. For example, if `adapter_weights=[0.5, 0.5]`, then the merged LoRA output is an average of both LoRAs. Try adjusting the adapter weights to see how it affects the generated image!
|
||||
|
||||
```py
|
||||
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
|
||||
@@ -47,19 +47,19 @@ image
|
||||
## add_weighted_adapter
|
||||
|
||||
> [!WARNING]
|
||||
> This is an experimental method that adds PEFTs [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method to Diffusers to enable more efficient merging methods. Check out this [issue](https://github.com/huggingface/diffusers/issues/6892) if you're interested in learning more about the motivation and design behind this integration.
|
||||
> This is an experimental method that adds PEFTs [`~peft.LoraModel.add_weighted_adapter`] method to Diffusers to enable more efficient merging methods. Check out this [issue](https://github.com/huggingface/diffusers/issues/6892) if you're interested in learning more about the motivation and design behind this integration.
|
||||
|
||||
The [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method provides access to more efficient merging method such as [TIES and DARE](https://huggingface.co/docs/peft/developer_guides/model_merging). To use these merging methods, make sure you have the latest stable version of Diffusers and PEFT installed.
|
||||
The [`~peft.LoraModel.add_weighted_adapter`] method provides access to more efficient merging method such as [TIES and DARE](https://huggingface.co/docs/peft/developer_guides/model_merging). To use these merging methods, make sure you have the latest stable version of Diffusers and PEFT installed.
|
||||
|
||||
```bash
|
||||
pip install -U diffusers peft
|
||||
```
|
||||
|
||||
There are three steps to merge LoRAs with the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method:
|
||||
There are three steps to merge LoRAs with the [`~peft.LoraModel.add_weighted_adapter`] method:
|
||||
|
||||
1. Create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the underlying model and LoRA checkpoint.
|
||||
1. Create a [`~peft.PeftModel`] from the underlying model and LoRA checkpoint.
|
||||
2. Load a base UNet model and the LoRA adapters.
|
||||
3. Merge the adapters using the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method and the merging method of your choice.
|
||||
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice.
|
||||
|
||||
Let's dive deeper into what these steps entail.
|
||||
|
||||
@@ -92,7 +92,7 @@ pipeline = DiffusionPipeline.from_pretrained(
|
||||
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
|
||||
```
|
||||
|
||||
Now you'll create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the loaded LoRA checkpoint by combining the SDXL UNet and the LoRA UNet from the pipeline.
|
||||
Now you'll create a [`~peft.PeftModel`] from the loaded LoRA checkpoint by combining the SDXL UNet and the LoRA UNet from the pipeline.
|
||||
|
||||
```python
|
||||
from peft import get_peft_model, LoraConfig
|
||||
@@ -112,7 +112,7 @@ ikea_peft_model.load_state_dict(original_state_dict, strict=True)
|
||||
> [!TIP]
|
||||
> You can optionally push the ikea_peft_model to the Hub by calling `ikea_peft_model.push_to_hub("ikea_peft_model", token=TOKEN)`.
|
||||
|
||||
Repeat this process to create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the [lordjia/by-feng-zikai](https://huggingface.co/lordjia/by-feng-zikai) LoRA.
|
||||
Repeat this process to create a [`~peft.PeftModel`] from the [lordjia/by-feng-zikai](https://huggingface.co/lordjia/by-feng-zikai) LoRA.
|
||||
|
||||
```python
|
||||
pipeline.delete_adapters("ikea")
|
||||
@@ -148,7 +148,7 @@ model = PeftModel.from_pretrained(base_unet, "stevhliu/ikea_peft_model", use_saf
|
||||
model.load_adapter("stevhliu/feng_peft_model", use_safetensors=True, subfolder="feng", adapter_name="feng")
|
||||
```
|
||||
|
||||
3. Merge the adapters using the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method and the merging method of your choice (learn more about other merging methods in this [blog post](https://huggingface.co/blog/peft_merging)). For this example, let's use the `"dare_linear"` method to merge the LoRAs.
|
||||
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice (learn more about other merging methods in this [blog post](https://huggingface.co/blog/peft_merging)). For this example, let's use the `"dare_linear"` method to merge the LoRAs.
|
||||
|
||||
> [!WARNING]
|
||||
> Keep in mind the LoRAs need to have the same rank to be merged!
|
||||
@@ -182,9 +182,9 @@ image
|
||||
|
||||
## fuse_lora
|
||||
|
||||
Both the [`~loaders.PeftAdapterMixin.set_adapters`] and [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
|
||||
Both the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
|
||||
|
||||
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
|
||||
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
|
||||
|
||||
For example, if you have a base model and adapters loaded and set as active with the following adapter weights:
|
||||
|
||||
@@ -199,7 +199,7 @@ pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_
|
||||
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
|
||||
```
|
||||
|
||||
Fuse these LoRAs into the UNet with the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method because it won’t work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
|
||||
Fuse these LoRAs into the UNet with the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method because it won’t work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
|
||||
|
||||
```py
|
||||
pipeline.fuse_lora(adapter_names=["ikea", "feng"], lora_scale=1.0)
|
||||
@@ -226,7 +226,7 @@ image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai"
|
||||
image
|
||||
```
|
||||
|
||||
You can call [`~~loaders.lora_base.LoraBaseMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
|
||||
You can call [`~loaders.StableDiffusionLoraLoaderMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
|
||||
|
||||
```py
|
||||
pipeline.unfuse_lora()
|
||||
|
||||
@@ -71,7 +71,7 @@ from diffusers.utils.import_utils import is_xformers_available
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -79,7 +79,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -71,7 +71,6 @@ Please also check out our [Community Scripts](https://github.com/huggingface/dif
|
||||
| Stable Diffusion BoxDiff Pipeline | Training-free controlled generation with bounding boxes using [BoxDiff](https://github.com/showlab/BoxDiff) | [Stable Diffusion BoxDiff Pipeline](#stable-diffusion-boxdiff) | - | [Jingyang Zhang](https://github.com/zjysteven/) |
|
||||
| FRESCO V2V Pipeline | Implementation of [[CVPR 2024] FRESCO: Spatial-Temporal Correspondence for Zero-Shot Video Translation](https://arxiv.org/abs/2403.12962) | [FRESCO V2V Pipeline](#fresco) | - | [Yifan Zhou](https://github.com/SingleZombie) |
|
||||
| AnimateDiff IPEX Pipeline | Accelerate AnimateDiff inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [AnimateDiff on IPEX](#animatediff-on-ipex) | - | [Dan Li](https://github.com/ustcuna/) |
|
||||
| HunyuanDiT Differential Diffusion Pipeline | Applies [Differential Diffsuion](https://github.com/exx8/differential-diffusion) to [HunyuanDiT](https://github.com/huggingface/diffusers/pull/8240). | [HunyuanDiT with Differential Diffusion](#hunyuandit-with-differential-diffusion) | [](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing) | [Monjoy Choudhury](https://github.com/MnCSSJ4x) |
|
||||
|
||||
To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly.
|
||||
|
||||
@@ -1647,6 +1646,7 @@ from diffusers import DiffusionPipeline
|
||||
scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1",
|
||||
subfolder="scheduler")
|
||||
|
||||
|
||||
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
|
||||
custom_pipeline="stable_diffusion_tensorrt_img2img",
|
||||
variant='fp16',
|
||||
@@ -1661,6 +1661,7 @@ pipe = pipe.to("cuda")
|
||||
url = "https://pajoca.com/wp-content/uploads/2022/09/tekito-yamakawa-1.png"
|
||||
response = requests.get(url)
|
||||
input_image = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
|
||||
prompt = "photorealistic new zealand hills"
|
||||
image = pipe(prompt, image=input_image, strength=0.75,).images[0]
|
||||
image.save('tensorrt_img2img_new_zealand_hills.png')
|
||||
@@ -4208,52 +4209,6 @@ print("Latency of AnimateDiffPipelineIpex--fp32", latency, "s for total", step,
|
||||
latency = elapsed_time(pipe4, num_inference_steps=step)
|
||||
print("Latency of AnimateDiffPipeline--fp32",latency, "s for total", step, "steps")
|
||||
```
|
||||
### HunyuanDiT with Differential Diffusion
|
||||
|
||||
#### Usage
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FlowMatchEulerDiscreteScheduler
|
||||
from diffusers.utils import load_image
|
||||
from PIL import Image
|
||||
from torchvision import transforms
|
||||
|
||||
from pipeline_hunyuandit_differential_img2img import (
|
||||
HunyuanDiTDifferentialImg2ImgPipeline,
|
||||
)
|
||||
|
||||
|
||||
pipe = HunyuanDiTDifferentialImg2ImgPipeline.from_pretrained(
|
||||
"Tencent-Hunyuan/HunyuanDiT-Diffusers", torch_dtype=torch.float16
|
||||
).to("cuda")
|
||||
|
||||
|
||||
source_image = load_image(
|
||||
"https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/20240329211129_4024911930.png"
|
||||
)
|
||||
map = load_image(
|
||||
"https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/gradient_mask_2.png"
|
||||
)
|
||||
prompt = "a green pear"
|
||||
negative_prompt = "blurry"
|
||||
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
image=source_image,
|
||||
num_inference_steps=28,
|
||||
guidance_scale=4.5,
|
||||
strength=1.0,
|
||||
map=map,
|
||||
).images[0]
|
||||
```
|
||||
|
||||
|  |  |  |
|
||||
| ------------------------------------------------------------------------------------------ | --------------------------------------------------------------------------------------- | ---------------------------------------------------------------------------------------- |
|
||||
| Gradient | Input | Output |
|
||||
|
||||
A colab notebook demonstrating all results can be found [here](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing). Depth Maps have also been added in the same colab.
|
||||
|
||||
# Perturbed-Attention Guidance
|
||||
|
||||
@@ -4330,4 +4285,4 @@ grid_image.save(grid_dir + "sample.png")
|
||||
|
||||
`pag_scale` : guidance scale of PAG (ex: 5.0)
|
||||
|
||||
`pag_applied_layers_index` : index of the layer to apply perturbation (ex: ['m0'])
|
||||
`pag_applied_layers_index` : index of the layer to apply perturbation (ex: ['m0'])
|
||||
@@ -43,8 +43,7 @@ from diffusers.utils import BaseOutput, check_min_version
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
|
||||
check_min_version("0.30.0")
|
||||
|
||||
class MarigoldDepthOutput(BaseOutput):
|
||||
"""
|
||||
|
||||
File diff suppressed because it is too large
Load Diff
@@ -73,7 +73,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -66,7 +66,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -79,7 +79,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -72,7 +72,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -78,7 +78,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -60,7 +60,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -60,7 +60,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = logging.getLogger(__name__)
|
||||
|
||||
|
||||
@@ -61,7 +61,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
if is_torch_npu_available():
|
||||
|
||||
@@ -63,7 +63,7 @@ from diffusers.utils.import_utils import is_xformers_available
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -1,195 +0,0 @@
|
||||
# DreamBooth training example for FLUX.1 [dev]
|
||||
|
||||
[DreamBooth](https://arxiv.org/abs/2208.12242) is a method to personalize text2image models like stable diffusion given just a few (3~5) images of a subject.
|
||||
|
||||
The `train_dreambooth_flux.py` script shows how to implement the training procedure and adapt it for [FLUX.1 [dev]](https://blackforestlabs.ai/announcing-black-forest-labs/). We also provide a LoRA implementation in the `train_dreambooth_lora_flux.py` script.
|
||||
> [!NOTE]
|
||||
> **Memory consumption**
|
||||
>
|
||||
> Flux can be quite expensive to run on consumer hardware devices and as a result finetuning it comes with high memory requirements -
|
||||
> a LoRA with a rank of 16 (w/ all components trained) can exceed 40GB of VRAM for training.
|
||||
> For more tips & guidance on training on a resource-constrained device please visit [`@bghira`'s guide](https://github.com/bghira/SimpleTuner/blob/main/documentation/quickstart/FLUX.md)
|
||||
|
||||
|
||||
> [!NOTE]
|
||||
> **Gated model**
|
||||
>
|
||||
> As the model is gated, before using it with diffusers you first need to go to the [FLUX.1 [dev] Hugging Face page](https://huggingface.co/black-forest-labs/FLUX.1-dev), fill in the form and accept the gate. Once you are in, you need to log in so that your system knows you’ve accepted the gate. Use the command below to log in:
|
||||
|
||||
```bash
|
||||
huggingface-cli login
|
||||
```
|
||||
|
||||
This will also allow us to push the trained model parameters to the Hugging Face Hub platform.
|
||||
|
||||
## Running locally with PyTorch
|
||||
|
||||
### Installing the dependencies
|
||||
|
||||
Before running the scripts, make sure to install the library's training dependencies:
|
||||
|
||||
**Important**
|
||||
|
||||
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
|
||||
|
||||
```bash
|
||||
git clone https://github.com/huggingface/diffusers
|
||||
cd diffusers
|
||||
pip install -e .
|
||||
```
|
||||
|
||||
Then cd in the `examples/dreambooth` folder and run
|
||||
```bash
|
||||
pip install -r requirements_flux.txt
|
||||
```
|
||||
|
||||
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
|
||||
|
||||
```bash
|
||||
accelerate config
|
||||
```
|
||||
|
||||
Or for a default accelerate configuration without answering questions about your environment
|
||||
|
||||
```bash
|
||||
accelerate config default
|
||||
```
|
||||
|
||||
Or if your environment doesn't support an interactive shell (e.g., a notebook)
|
||||
|
||||
```python
|
||||
from accelerate.utils import write_basic_config
|
||||
write_basic_config()
|
||||
```
|
||||
|
||||
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
|
||||
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
|
||||
|
||||
|
||||
### Dog toy example
|
||||
|
||||
Now let's get our dataset. For this example we will use some dog images: https://huggingface.co/datasets/diffusers/dog-example.
|
||||
|
||||
Let's first download it locally:
|
||||
|
||||
```python
|
||||
from huggingface_hub import snapshot_download
|
||||
|
||||
local_dir = "./dog"
|
||||
snapshot_download(
|
||||
"diffusers/dog-example",
|
||||
local_dir=local_dir, repo_type="dataset",
|
||||
ignore_patterns=".gitattributes",
|
||||
)
|
||||
```
|
||||
|
||||
This will also allow us to push the trained LoRA parameters to the Hugging Face Hub platform.
|
||||
|
||||
Now, we can launch training using:
|
||||
|
||||
```bash
|
||||
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
|
||||
export INSTANCE_DIR="dog"
|
||||
export OUTPUT_DIR="trained-flux"
|
||||
|
||||
accelerate launch train_dreambooth_flux.py \
|
||||
--pretrained_model_name_or_path=$MODEL_NAME \
|
||||
--instance_data_dir=$INSTANCE_DIR \
|
||||
--output_dir=$OUTPUT_DIR \
|
||||
--mixed_precision="bf16" \
|
||||
--instance_prompt="a photo of sks dog" \
|
||||
--resolution=1024 \
|
||||
--train_batch_size=1 \
|
||||
--gradient_accumulation_steps=4 \
|
||||
--learning_rate=1e-4 \
|
||||
--report_to="wandb" \
|
||||
--lr_scheduler="constant" \
|
||||
--lr_warmup_steps=0 \
|
||||
--max_train_steps=500 \
|
||||
--validation_prompt="A photo of sks dog in a bucket" \
|
||||
--validation_epochs=25 \
|
||||
--seed="0" \
|
||||
--push_to_hub
|
||||
```
|
||||
|
||||
To better track our training experiments, we're using the following flags in the command above:
|
||||
|
||||
* `report_to="wandb` will ensure the training runs are tracked on Weights and Biases. To use it, be sure to install `wandb` with `pip install wandb`.
|
||||
* `validation_prompt` and `validation_epochs` to allow the script to do a few validation inference runs. This allows us to qualitatively check if the training is progressing as expected.
|
||||
|
||||
> [!NOTE]
|
||||
> If you want to train using long prompts with the T5 text encoder, you can use `--max_sequence_length` to set the token limit. The default is 77, but it can be increased to as high as 512. Note that this will use more resources and may slow down the training in some cases.
|
||||
|
||||
> [!TIP]
|
||||
> You can pass `--use_8bit_adam` to reduce the memory requirements of training. Make sure to install `bitsandbytes` if you want to do so.
|
||||
|
||||
## LoRA + DreamBooth
|
||||
|
||||
[LoRA](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) is a popular parameter-efficient fine-tuning technique that allows you to achieve full-finetuning like performance but with a fraction of learnable parameters.
|
||||
|
||||
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
|
||||
|
||||
To perform DreamBooth with LoRA, run:
|
||||
|
||||
```bash
|
||||
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
|
||||
export INSTANCE_DIR="dog"
|
||||
export OUTPUT_DIR="trained-flux-lora"
|
||||
|
||||
accelerate launch train_dreambooth_lora_flux.py \
|
||||
--pretrained_model_name_or_path=$MODEL_NAME \
|
||||
--instance_data_dir=$INSTANCE_DIR \
|
||||
--output_dir=$OUTPUT_DIR \
|
||||
--mixed_precision="bf16" \
|
||||
--instance_prompt="a photo of sks dog" \
|
||||
--resolution=512 \
|
||||
--train_batch_size=1 \
|
||||
--gradient_accumulation_steps=4 \
|
||||
--learning_rate=1e-5 \
|
||||
--report_to="wandb" \
|
||||
--lr_scheduler="constant" \
|
||||
--lr_warmup_steps=0 \
|
||||
--max_train_steps=500 \
|
||||
--validation_prompt="A photo of sks dog in a bucket" \
|
||||
--validation_epochs=25 \
|
||||
--seed="0" \
|
||||
--push_to_hub
|
||||
```
|
||||
|
||||
### Text Encoder Training
|
||||
|
||||
Alongside the transformer, fine-tuning of the CLIP text encoder is also supported.
|
||||
To do so, just specify `--train_text_encoder` while launching training. Please keep the following points in mind:
|
||||
|
||||
> [!NOTE]
|
||||
> FLUX.1 has 2 text encoders (CLIP L/14 and T5-v1.1-XXL).
|
||||
By enabling `--train_text_encoder`, fine-tuning of the **CLIP encoder** is performed.
|
||||
> At the moment, T5 fine-tuning is not supported and weights remain frozen when text encoder training is enabled.
|
||||
|
||||
To perform DreamBooth LoRA with text-encoder training, run:
|
||||
```bash
|
||||
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
|
||||
export OUTPUT_DIR="trained-flux-dev-dreambooth-lora"
|
||||
|
||||
accelerate launch train_dreambooth_lora_flux.py \
|
||||
--pretrained_model_name_or_path=$MODEL_NAME \
|
||||
--instance_data_dir=$INSTANCE_DIR \
|
||||
--output_dir=$OUTPUT_DIR \
|
||||
--mixed_precision="bf16" \
|
||||
--train_text_encoder\
|
||||
--instance_prompt="a photo of sks dog" \
|
||||
--resolution=512 \
|
||||
--train_batch_size=1 \
|
||||
--gradient_accumulation_steps=4 \
|
||||
--learning_rate=1e-5 \
|
||||
--report_to="wandb" \
|
||||
--lr_scheduler="constant" \
|
||||
--lr_warmup_steps=0 \
|
||||
--max_train_steps=500 \
|
||||
--validation_prompt="A photo of sks dog in a bucket" \
|
||||
--seed="0" \
|
||||
--push_to_hub
|
||||
```
|
||||
|
||||
## Other notes
|
||||
Thanks to `bghira` for their help with reviewing & insight sharing ♥️
|
||||
@@ -1,8 +0,0 @@
|
||||
accelerate>=0.31.0
|
||||
torchvision
|
||||
transformers>=4.41.2
|
||||
ftfy
|
||||
tensorboard
|
||||
Jinja2
|
||||
peft>=0.11.1
|
||||
sentencepiece
|
||||
@@ -1,203 +0,0 @@
|
||||
# coding=utf-8
|
||||
# Copyright 2024 HuggingFace Inc.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
import logging
|
||||
import os
|
||||
import shutil
|
||||
import sys
|
||||
import tempfile
|
||||
|
||||
from diffusers import DiffusionPipeline, FluxTransformer2DModel
|
||||
|
||||
|
||||
sys.path.append("..")
|
||||
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
|
||||
|
||||
|
||||
logging.basicConfig(level=logging.DEBUG)
|
||||
|
||||
logger = logging.getLogger()
|
||||
stream_handler = logging.StreamHandler(sys.stdout)
|
||||
logger.addHandler(stream_handler)
|
||||
|
||||
|
||||
class DreamBoothFlux(ExamplesTestsAccelerate):
|
||||
instance_data_dir = "docs/source/en/imgs"
|
||||
instance_prompt = "photo"
|
||||
pretrained_model_name_or_path = "hf-internal-testing/tiny-flux-pipe"
|
||||
script_path = "examples/dreambooth/train_dreambooth_flux.py"
|
||||
|
||||
def test_dreambooth(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
|
||||
--instance_data_dir {self.instance_data_dir}
|
||||
--instance_prompt {self.instance_prompt}
|
||||
--resolution 64
|
||||
--train_batch_size 1
|
||||
--gradient_accumulation_steps 1
|
||||
--max_train_steps 2
|
||||
--learning_rate 5.0e-04
|
||||
--scale_lr
|
||||
--lr_scheduler constant
|
||||
--lr_warmup_steps 0
|
||||
--output_dir {tmpdir}
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
# save_pretrained smoke test
|
||||
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "transformer", "diffusion_pytorch_model.safetensors")))
|
||||
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "scheduler", "scheduler_config.json")))
|
||||
|
||||
def test_dreambooth_checkpointing(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
# Run training script with checkpointing
|
||||
# max_train_steps == 4, checkpointing_steps == 2
|
||||
# Should create checkpoints at steps 2, 4
|
||||
|
||||
initial_run_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
|
||||
--instance_data_dir {self.instance_data_dir}
|
||||
--instance_prompt {self.instance_prompt}
|
||||
--resolution 64
|
||||
--train_batch_size 1
|
||||
--gradient_accumulation_steps 1
|
||||
--max_train_steps 4
|
||||
--learning_rate 5.0e-04
|
||||
--scale_lr
|
||||
--lr_scheduler constant
|
||||
--lr_warmup_steps 0
|
||||
--output_dir {tmpdir}
|
||||
--checkpointing_steps=2
|
||||
--seed=0
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + initial_run_args)
|
||||
|
||||
# check can run the original fully trained output pipeline
|
||||
pipe = DiffusionPipeline.from_pretrained(tmpdir)
|
||||
pipe(self.instance_prompt, num_inference_steps=1)
|
||||
|
||||
# check checkpoint directories exist
|
||||
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
|
||||
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
|
||||
|
||||
# check can run an intermediate checkpoint
|
||||
transformer = FluxTransformer2DModel.from_pretrained(tmpdir, subfolder="checkpoint-2/transformer")
|
||||
pipe = DiffusionPipeline.from_pretrained(self.pretrained_model_name_or_path, transformer=transformer)
|
||||
pipe(self.instance_prompt, num_inference_steps=1)
|
||||
|
||||
# Remove checkpoint 2 so that we can check only later checkpoints exist after resuming
|
||||
shutil.rmtree(os.path.join(tmpdir, "checkpoint-2"))
|
||||
|
||||
# Run training script for 7 total steps resuming from checkpoint 4
|
||||
|
||||
resume_run_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
|
||||
--instance_data_dir {self.instance_data_dir}
|
||||
--instance_prompt {self.instance_prompt}
|
||||
--resolution 64
|
||||
--train_batch_size 1
|
||||
--gradient_accumulation_steps 1
|
||||
--max_train_steps 6
|
||||
--learning_rate 5.0e-04
|
||||
--scale_lr
|
||||
--lr_scheduler constant
|
||||
--lr_warmup_steps 0
|
||||
--output_dir {tmpdir}
|
||||
--checkpointing_steps=2
|
||||
--resume_from_checkpoint=checkpoint-4
|
||||
--seed=0
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + resume_run_args)
|
||||
|
||||
# check can run new fully trained pipeline
|
||||
pipe = DiffusionPipeline.from_pretrained(tmpdir)
|
||||
pipe(self.instance_prompt, num_inference_steps=1)
|
||||
|
||||
# check old checkpoints do not exist
|
||||
self.assertFalse(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
|
||||
|
||||
# check new checkpoints exist
|
||||
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
|
||||
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-6")))
|
||||
|
||||
def test_dreambooth_checkpointing_checkpoints_total_limit(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=6
|
||||
--checkpoints_total_limit=2
|
||||
--checkpointing_steps=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
|
||||
self.assertEqual(
|
||||
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
|
||||
{"checkpoint-4", "checkpoint-6"},
|
||||
)
|
||||
|
||||
def test_dreambooth_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=4
|
||||
--checkpointing_steps=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
|
||||
self.assertEqual(
|
||||
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
|
||||
{"checkpoint-2", "checkpoint-4"},
|
||||
)
|
||||
|
||||
resume_run_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=8
|
||||
--checkpointing_steps=2
|
||||
--resume_from_checkpoint=checkpoint-4
|
||||
--checkpoints_total_limit=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + resume_run_args)
|
||||
|
||||
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-6", "checkpoint-8"})
|
||||
@@ -1,165 +0,0 @@
|
||||
# coding=utf-8
|
||||
# Copyright 2024 HuggingFace Inc.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
import logging
|
||||
import os
|
||||
import sys
|
||||
import tempfile
|
||||
|
||||
import safetensors
|
||||
|
||||
|
||||
sys.path.append("..")
|
||||
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
|
||||
|
||||
|
||||
logging.basicConfig(level=logging.DEBUG)
|
||||
|
||||
logger = logging.getLogger()
|
||||
stream_handler = logging.StreamHandler(sys.stdout)
|
||||
logger.addHandler(stream_handler)
|
||||
|
||||
|
||||
class DreamBoothLoRAFlux(ExamplesTestsAccelerate):
|
||||
instance_data_dir = "docs/source/en/imgs"
|
||||
instance_prompt = "photo"
|
||||
pretrained_model_name_or_path = "hf-internal-testing/tiny-flux-pipe"
|
||||
script_path = "examples/dreambooth/train_dreambooth_lora_flux.py"
|
||||
|
||||
def test_dreambooth_lora_flux(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
|
||||
--instance_data_dir {self.instance_data_dir}
|
||||
--instance_prompt {self.instance_prompt}
|
||||
--resolution 64
|
||||
--train_batch_size 1
|
||||
--gradient_accumulation_steps 1
|
||||
--max_train_steps 2
|
||||
--learning_rate 5.0e-04
|
||||
--scale_lr
|
||||
--lr_scheduler constant
|
||||
--lr_warmup_steps 0
|
||||
--output_dir {tmpdir}
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
# save_pretrained smoke test
|
||||
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
|
||||
|
||||
# make sure the state_dict has the correct naming in the parameters.
|
||||
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
|
||||
is_lora = all("lora" in k for k in lora_state_dict.keys())
|
||||
self.assertTrue(is_lora)
|
||||
|
||||
# when not training the text encoder, all the parameters in the state dict should start
|
||||
# with `"transformer"` in their names.
|
||||
starts_with_transformer = all(key.startswith("transformer") for key in lora_state_dict.keys())
|
||||
self.assertTrue(starts_with_transformer)
|
||||
|
||||
def test_dreambooth_lora_text_encoder_flux(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
|
||||
--instance_data_dir {self.instance_data_dir}
|
||||
--instance_prompt {self.instance_prompt}
|
||||
--resolution 64
|
||||
--train_batch_size 1
|
||||
--train_text_encoder
|
||||
--gradient_accumulation_steps 1
|
||||
--max_train_steps 2
|
||||
--learning_rate 5.0e-04
|
||||
--scale_lr
|
||||
--lr_scheduler constant
|
||||
--lr_warmup_steps 0
|
||||
--output_dir {tmpdir}
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
# save_pretrained smoke test
|
||||
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
|
||||
|
||||
# make sure the state_dict has the correct naming in the parameters.
|
||||
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
|
||||
is_lora = all("lora" in k for k in lora_state_dict.keys())
|
||||
self.assertTrue(is_lora)
|
||||
|
||||
starts_with_expected_prefix = all(
|
||||
(key.startswith("transformer") or key.startswith("text_encoder")) for key in lora_state_dict.keys()
|
||||
)
|
||||
self.assertTrue(starts_with_expected_prefix)
|
||||
|
||||
def test_dreambooth_lora_flux_checkpointing_checkpoints_total_limit(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=6
|
||||
--checkpoints_total_limit=2
|
||||
--checkpointing_steps=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
|
||||
self.assertEqual(
|
||||
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
|
||||
{"checkpoint-4", "checkpoint-6"},
|
||||
)
|
||||
|
||||
def test_dreambooth_lora_flux_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
|
||||
with tempfile.TemporaryDirectory() as tmpdir:
|
||||
test_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=4
|
||||
--checkpointing_steps=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + test_args)
|
||||
|
||||
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-2", "checkpoint-4"})
|
||||
|
||||
resume_run_args = f"""
|
||||
{self.script_path}
|
||||
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
|
||||
--instance_data_dir={self.instance_data_dir}
|
||||
--output_dir={tmpdir}
|
||||
--instance_prompt={self.instance_prompt}
|
||||
--resolution=64
|
||||
--train_batch_size=1
|
||||
--gradient_accumulation_steps=1
|
||||
--max_train_steps=8
|
||||
--checkpointing_steps=2
|
||||
--resume_from_checkpoint=checkpoint-4
|
||||
--checkpoints_total_limit=2
|
||||
""".split()
|
||||
|
||||
run_command(self._launch_args + resume_run_args)
|
||||
|
||||
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-6", "checkpoint-8"})
|
||||
@@ -63,7 +63,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -35,7 +35,7 @@ from diffusers.utils import check_min_version
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
# Cache compiled models across invocations of this script.
|
||||
cc.initialize_cache(os.path.expanduser("~/.cache/jax/compilation_cache"))
|
||||
|
||||
File diff suppressed because it is too large
Load Diff
@@ -70,7 +70,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
File diff suppressed because it is too large
Load Diff
@@ -72,7 +72,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -78,7 +78,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -64,7 +64,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -57,7 +57,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -60,7 +60,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -52,7 +52,7 @@ if is_wandb_available():
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -51,7 +51,7 @@ if is_wandb_available():
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -60,7 +60,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -57,7 +57,7 @@ if is_wandb_available():
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
@@ -826,22 +826,17 @@ def main():
|
||||
)
|
||||
|
||||
# Scheduler and math around the number of training steps.
|
||||
# Check the PR https://github.com/huggingface/diffusers/pull/8312 for detailed explanation.
|
||||
num_warmup_steps_for_scheduler = args.lr_warmup_steps * accelerator.num_processes
|
||||
overrode_max_train_steps = False
|
||||
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
|
||||
if args.max_train_steps is None:
|
||||
len_train_dataloader_after_sharding = math.ceil(len(train_dataloader) / accelerator.num_processes)
|
||||
num_update_steps_per_epoch = math.ceil(len_train_dataloader_after_sharding / args.gradient_accumulation_steps)
|
||||
num_training_steps_for_scheduler = (
|
||||
args.num_train_epochs * num_update_steps_per_epoch * accelerator.num_processes
|
||||
)
|
||||
else:
|
||||
num_training_steps_for_scheduler = args.max_train_steps * accelerator.num_processes
|
||||
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
|
||||
overrode_max_train_steps = True
|
||||
|
||||
lr_scheduler = get_scheduler(
|
||||
args.lr_scheduler,
|
||||
optimizer=optimizer,
|
||||
num_warmup_steps=num_warmup_steps_for_scheduler,
|
||||
num_training_steps=num_training_steps_for_scheduler,
|
||||
num_warmup_steps=args.lr_warmup_steps * accelerator.num_processes,
|
||||
num_training_steps=args.max_train_steps * accelerator.num_processes,
|
||||
)
|
||||
|
||||
# Prepare everything with our `accelerator`.
|
||||
@@ -871,14 +866,8 @@ def main():
|
||||
|
||||
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
|
||||
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
|
||||
if args.max_train_steps is None:
|
||||
if overrode_max_train_steps:
|
||||
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
|
||||
if num_training_steps_for_scheduler != args.max_train_steps * accelerator.num_processes:
|
||||
logger.warning(
|
||||
f"The length of the 'train_dataloader' after 'accelerator.prepare' ({len(train_dataloader)}) does not match "
|
||||
f"the expected length ({len_train_dataloader_after_sharding}) when the learning rate scheduler was created. "
|
||||
f"This inconsistency may result in the learning rate scheduler not functioning properly."
|
||||
)
|
||||
# Afterwards we recalculate our number of training epochs
|
||||
args.num_train_epochs = math.ceil(args.max_train_steps / num_update_steps_per_epoch)
|
||||
|
||||
|
||||
@@ -49,7 +49,7 @@ from diffusers.utils import check_min_version
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = logging.getLogger(__name__)
|
||||
|
||||
|
||||
@@ -56,7 +56,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -68,7 +68,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
if is_torch_npu_available():
|
||||
@@ -478,7 +478,7 @@ def parse_args(input_args=None):
|
||||
parser.add_argument(
|
||||
"--debug_loss",
|
||||
action="store_true",
|
||||
help="debug loss for each image, if filenames are available in the dataset",
|
||||
help="debug loss for each image, if filenames are awailable in the dataset",
|
||||
)
|
||||
|
||||
if input_args is not None:
|
||||
|
||||
@@ -55,7 +55,7 @@ from diffusers.utils.torch_utils import is_compiled_module
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
if is_torch_npu_available():
|
||||
|
||||
@@ -23,25 +23,4 @@ accelerate launch textual_inversion_sdxl.py \
|
||||
--output_dir="./textual_inversion_cat_sdxl"
|
||||
```
|
||||
|
||||
Training of both text encoders is supported.
|
||||
|
||||
### Inference Example
|
||||
|
||||
Once you have trained a model using above command, the inference can be done simply using the `StableDiffusionXLPipeline`.
|
||||
Make sure to include the `placeholder_token` in your prompt.
|
||||
|
||||
```python
|
||||
from diffusers import StableDiffusionXLPipeline
|
||||
import torch
|
||||
|
||||
model_id = "./textual_inversion_cat_sdxl"
|
||||
pipe = StableDiffusionXLPipeline.from_pretrained(model_id,torch_dtype=torch.float16).to("cuda")
|
||||
|
||||
prompt = "A <cat-toy> backpack"
|
||||
|
||||
image = pipe(prompt, num_inference_steps=50, guidance_scale=7.5).images[0]
|
||||
image.save("cat-backpack.png")
|
||||
|
||||
image = pipe(prompt="", prompt_2=prompt, num_inference_steps=50, guidance_scale=7.5).images[0]
|
||||
image.save("cat-backpack-prompt_2.png")
|
||||
```
|
||||
For now, only training of the first text encoder is supported.
|
||||
@@ -81,7 +81,7 @@ else:
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
@@ -56,7 +56,7 @@ else:
|
||||
# ------------------------------------------------------------------------------
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = logging.getLogger(__name__)
|
||||
|
||||
|
||||
@@ -76,7 +76,7 @@ else:
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
@@ -135,7 +135,7 @@ def log_validation(
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
text_encoder=accelerator.unwrap_model(text_encoder_1),
|
||||
text_encoder_2=accelerator.unwrap_model(text_encoder_2),
|
||||
text_encoder_2=text_encoder_2,
|
||||
tokenizer=tokenizer_1,
|
||||
tokenizer_2=tokenizer_2,
|
||||
unet=unet,
|
||||
@@ -678,54 +678,36 @@ def main():
|
||||
f"The tokenizer already contains the token {args.placeholder_token}. Please pass a different"
|
||||
" `placeholder_token` that is not already in the tokenizer."
|
||||
)
|
||||
num_added_tokens = tokenizer_2.add_tokens(placeholder_tokens)
|
||||
if num_added_tokens != args.num_vectors:
|
||||
raise ValueError(
|
||||
f"The 2nd tokenizer already contains the token {args.placeholder_token}. Please pass a different"
|
||||
" `placeholder_token` that is not already in the tokenizer."
|
||||
)
|
||||
|
||||
# Convert the initializer_token, placeholder_token to ids
|
||||
token_ids = tokenizer_1.encode(args.initializer_token, add_special_tokens=False)
|
||||
token_ids_2 = tokenizer_2.encode(args.initializer_token, add_special_tokens=False)
|
||||
|
||||
# Check if initializer_token is a single token or a sequence of tokens
|
||||
if len(token_ids) > 1 or len(token_ids_2) > 1:
|
||||
if len(token_ids) > 1:
|
||||
raise ValueError("The initializer token must be a single token.")
|
||||
|
||||
initializer_token_id = token_ids[0]
|
||||
placeholder_token_ids = tokenizer_1.convert_tokens_to_ids(placeholder_tokens)
|
||||
initializer_token_id_2 = token_ids_2[0]
|
||||
placeholder_token_ids_2 = tokenizer_2.convert_tokens_to_ids(placeholder_tokens)
|
||||
|
||||
# Resize the token embeddings as we are adding new special tokens to the tokenizer
|
||||
text_encoder_1.resize_token_embeddings(len(tokenizer_1))
|
||||
text_encoder_2.resize_token_embeddings(len(tokenizer_2))
|
||||
|
||||
# Initialise the newly added placeholder token with the embeddings of the initializer token
|
||||
token_embeds = text_encoder_1.get_input_embeddings().weight.data
|
||||
token_embeds_2 = text_encoder_2.get_input_embeddings().weight.data
|
||||
with torch.no_grad():
|
||||
for token_id in placeholder_token_ids:
|
||||
token_embeds[token_id] = token_embeds[initializer_token_id].clone()
|
||||
for token_id in placeholder_token_ids_2:
|
||||
token_embeds_2[token_id] = token_embeds_2[initializer_token_id_2].clone()
|
||||
|
||||
# Freeze vae and unet
|
||||
vae.requires_grad_(False)
|
||||
unet.requires_grad_(False)
|
||||
|
||||
text_encoder_2.requires_grad_(False)
|
||||
# Freeze all parameters except for the token embeddings in text encoder
|
||||
text_encoder_1.text_model.encoder.requires_grad_(False)
|
||||
text_encoder_1.text_model.final_layer_norm.requires_grad_(False)
|
||||
text_encoder_1.text_model.embeddings.position_embedding.requires_grad_(False)
|
||||
text_encoder_2.text_model.encoder.requires_grad_(False)
|
||||
text_encoder_2.text_model.final_layer_norm.requires_grad_(False)
|
||||
text_encoder_2.text_model.embeddings.position_embedding.requires_grad_(False)
|
||||
|
||||
if args.gradient_checkpointing:
|
||||
text_encoder_1.gradient_checkpointing_enable()
|
||||
text_encoder_2.gradient_checkpointing_enable()
|
||||
|
||||
if args.enable_xformers_memory_efficient_attention:
|
||||
if is_xformers_available():
|
||||
@@ -764,11 +746,7 @@ def main():
|
||||
optimizer_class = torch.optim.AdamW
|
||||
|
||||
optimizer = optimizer_class(
|
||||
# only optimize the embeddings
|
||||
[
|
||||
text_encoder_1.text_model.embeddings.token_embedding.weight,
|
||||
text_encoder_2.text_model.embeddings.token_embedding.weight,
|
||||
],
|
||||
text_encoder_1.get_input_embeddings().parameters(), # only optimize the embeddings
|
||||
lr=args.learning_rate,
|
||||
betas=(args.adam_beta1, args.adam_beta2),
|
||||
weight_decay=args.adam_weight_decay,
|
||||
@@ -808,10 +786,9 @@ def main():
|
||||
)
|
||||
|
||||
text_encoder_1.train()
|
||||
text_encoder_2.train()
|
||||
# Prepare everything with our `accelerator`.
|
||||
text_encoder_1, text_encoder_2, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
|
||||
text_encoder_1, text_encoder_2, optimizer, train_dataloader, lr_scheduler
|
||||
text_encoder_1, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
|
||||
text_encoder_1, optimizer, train_dataloader, lr_scheduler
|
||||
)
|
||||
|
||||
# For mixed precision training we cast all non-trainable weigths (vae, non-lora text_encoder and non-lora unet) to half-precision
|
||||
@@ -889,13 +866,11 @@ def main():
|
||||
|
||||
# keep original embeddings as reference
|
||||
orig_embeds_params = accelerator.unwrap_model(text_encoder_1).get_input_embeddings().weight.data.clone()
|
||||
orig_embeds_params_2 = accelerator.unwrap_model(text_encoder_2).get_input_embeddings().weight.data.clone()
|
||||
|
||||
for epoch in range(first_epoch, args.num_train_epochs):
|
||||
text_encoder_1.train()
|
||||
text_encoder_2.train()
|
||||
for step, batch in enumerate(train_dataloader):
|
||||
with accelerator.accumulate([text_encoder_1, text_encoder_2]):
|
||||
with accelerator.accumulate(text_encoder_1):
|
||||
# Convert images to latent space
|
||||
latents = vae.encode(batch["pixel_values"].to(dtype=weight_dtype)).latent_dist.sample().detach()
|
||||
latents = latents * vae.config.scaling_factor
|
||||
@@ -917,7 +892,9 @@ def main():
|
||||
.hidden_states[-2]
|
||||
.to(dtype=weight_dtype)
|
||||
)
|
||||
encoder_output_2 = text_encoder_2(batch["input_ids_2"], output_hidden_states=True)
|
||||
encoder_output_2 = text_encoder_2(
|
||||
batch["input_ids_2"].reshape(batch["input_ids_1"].shape[0], -1), output_hidden_states=True
|
||||
)
|
||||
encoder_hidden_states_2 = encoder_output_2.hidden_states[-2].to(dtype=weight_dtype)
|
||||
original_size = [
|
||||
(batch["original_size"][0][i].item(), batch["original_size"][1][i].item())
|
||||
@@ -961,16 +938,11 @@ def main():
|
||||
# Let's make sure we don't update any embedding weights besides the newly added token
|
||||
index_no_updates = torch.ones((len(tokenizer_1),), dtype=torch.bool)
|
||||
index_no_updates[min(placeholder_token_ids) : max(placeholder_token_ids) + 1] = False
|
||||
index_no_updates_2 = torch.ones((len(tokenizer_2),), dtype=torch.bool)
|
||||
index_no_updates_2[min(placeholder_token_ids_2) : max(placeholder_token_ids_2) + 1] = False
|
||||
|
||||
with torch.no_grad():
|
||||
accelerator.unwrap_model(text_encoder_1).get_input_embeddings().weight[
|
||||
index_no_updates
|
||||
] = orig_embeds_params[index_no_updates]
|
||||
accelerator.unwrap_model(text_encoder_2).get_input_embeddings().weight[
|
||||
index_no_updates_2
|
||||
] = orig_embeds_params_2[index_no_updates_2]
|
||||
|
||||
# Checks if the accelerator has performed an optimization step behind the scenes
|
||||
if accelerator.sync_gradients:
|
||||
@@ -988,16 +960,6 @@ def main():
|
||||
save_path,
|
||||
safe_serialization=True,
|
||||
)
|
||||
weight_name = f"learned_embeds_2-steps-{global_step}.safetensors"
|
||||
save_path = os.path.join(args.output_dir, weight_name)
|
||||
save_progress(
|
||||
text_encoder_2,
|
||||
placeholder_token_ids_2,
|
||||
accelerator,
|
||||
args,
|
||||
save_path,
|
||||
safe_serialization=True,
|
||||
)
|
||||
|
||||
if accelerator.is_main_process:
|
||||
if global_step % args.checkpointing_steps == 0:
|
||||
@@ -1072,7 +1034,7 @@ def main():
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
text_encoder=accelerator.unwrap_model(text_encoder_1),
|
||||
text_encoder_2=accelerator.unwrap_model(text_encoder_2),
|
||||
text_encoder_2=text_encoder_2,
|
||||
vae=vae,
|
||||
unet=unet,
|
||||
tokenizer=tokenizer_1,
|
||||
@@ -1090,16 +1052,6 @@ def main():
|
||||
save_path,
|
||||
safe_serialization=True,
|
||||
)
|
||||
weight_name = "learned_embeds_2.safetensors"
|
||||
save_path = os.path.join(args.output_dir, weight_name)
|
||||
save_progress(
|
||||
text_encoder_2,
|
||||
placeholder_token_ids_2,
|
||||
accelerator,
|
||||
args,
|
||||
save_path,
|
||||
safe_serialization=True,
|
||||
)
|
||||
|
||||
if args.push_to_hub:
|
||||
save_model_card(
|
||||
|
||||
@@ -29,7 +29,7 @@ from diffusers.utils.import_utils import is_xformers_available
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -50,7 +50,7 @@ if is_wandb_available():
|
||||
import wandb
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -50,7 +50,7 @@ if is_wandb_available():
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -51,7 +51,7 @@ if is_wandb_available():
|
||||
|
||||
|
||||
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
|
||||
check_min_version("0.30.0.dev0")
|
||||
check_min_version("0.30.0")
|
||||
|
||||
logger = get_logger(__name__, log_level="INFO")
|
||||
|
||||
|
||||
@@ -86,6 +86,9 @@ TRANSFORMER_SPECIAL_KEYS_REMAP = {
|
||||
"key_layernorm_list": reassign_query_key_layernorm_inplace,
|
||||
"adaln_layer.adaLN_modulations": reassign_adaln_norm_inplace,
|
||||
"embed_tokens": remove_keys_inplace,
|
||||
"freqs_sin": remove_keys_inplace,
|
||||
"freqs_cos": remove_keys_inplace,
|
||||
"position_embedding": remove_keys_inplace,
|
||||
}
|
||||
|
||||
VAE_KEYS_RENAME_DICT = {
|
||||
@@ -123,11 +126,21 @@ def update_state_dict_inplace(state_dict: Dict[str, Any], old_key: str, new_key:
|
||||
state_dict[new_key] = state_dict.pop(old_key)
|
||||
|
||||
|
||||
def convert_transformer(ckpt_path: str):
|
||||
def convert_transformer(
|
||||
ckpt_path: str,
|
||||
num_layers: int,
|
||||
num_attention_heads: int,
|
||||
use_rotary_positional_embeddings: bool,
|
||||
dtype: torch.dtype,
|
||||
):
|
||||
PREFIX_KEY = "model.diffusion_model."
|
||||
|
||||
original_state_dict = get_state_dict(torch.load(ckpt_path, map_location="cpu", mmap=True))
|
||||
transformer = CogVideoXTransformer3DModel()
|
||||
transformer = CogVideoXTransformer3DModel(
|
||||
num_layers=num_layers,
|
||||
num_attention_heads=num_attention_heads,
|
||||
use_rotary_positional_embeddings=use_rotary_positional_embeddings,
|
||||
).to(dtype=dtype)
|
||||
|
||||
for key in list(original_state_dict.keys()):
|
||||
new_key = key[len(PREFIX_KEY) :]
|
||||
@@ -145,9 +158,9 @@ def convert_transformer(ckpt_path: str):
|
||||
return transformer
|
||||
|
||||
|
||||
def convert_vae(ckpt_path: str):
|
||||
def convert_vae(ckpt_path: str, scaling_factor: float, dtype: torch.dtype):
|
||||
original_state_dict = get_state_dict(torch.load(ckpt_path, map_location="cpu", mmap=True))
|
||||
vae = AutoencoderKLCogVideoX()
|
||||
vae = AutoencoderKLCogVideoX(scaling_factor=scaling_factor).to(dtype=dtype)
|
||||
|
||||
for key in list(original_state_dict.keys()):
|
||||
new_key = key[:]
|
||||
@@ -172,13 +185,26 @@ def get_args():
|
||||
)
|
||||
parser.add_argument("--vae_ckpt_path", type=str, default=None, help="Path to original vae checkpoint")
|
||||
parser.add_argument("--output_path", type=str, required=True, help="Path where converted model should be saved")
|
||||
parser.add_argument("--fp16", action="store_true", default=True, help="Whether to save the model weights in fp16")
|
||||
parser.add_argument("--fp16", action="store_true", default=False, help="Whether to save the model weights in fp16")
|
||||
parser.add_argument("--bf16", action="store_true", default=False, help="Whether to save the model weights in bf16")
|
||||
parser.add_argument(
|
||||
"--push_to_hub", action="store_true", default=False, help="Whether to push to HF Hub after saving"
|
||||
)
|
||||
parser.add_argument(
|
||||
"--text_encoder_cache_dir", type=str, default=None, help="Path to text encoder cache directory"
|
||||
)
|
||||
# For CogVideoX-2B, num_layers is 30. For 5B, it is 42
|
||||
parser.add_argument("--num_layers", type=int, default=30, help="Number of transformer blocks")
|
||||
# For CogVideoX-2B, num_attention_heads is 30. For 5B, it is 48
|
||||
parser.add_argument("--num_attention_heads", type=int, default=30, help="Number of attention heads")
|
||||
# For CogVideoX-2B, use_rotary_positional_embeddings is False. For 5B, it is True
|
||||
parser.add_argument(
|
||||
"--use_rotary_positional_embeddings", action="store_true", default=False, help="Whether to use RoPE or not"
|
||||
)
|
||||
# For CogVideoX-2B, scaling_factor is 1.15258426. For 5B, it is 0.7
|
||||
parser.add_argument("--scaling_factor", type=float, default=1.15258426, help="Scaling factor in the VAE")
|
||||
# For CogVideoX-2B, snr_shift_scale is 3.0. For 5B, it is 1.0
|
||||
parser.add_argument("--snr_shift_scale", type=float, default=3.0, help="Scaling factor in the VAE")
|
||||
return parser.parse_args()
|
||||
|
||||
|
||||
@@ -188,18 +214,33 @@ if __name__ == "__main__":
|
||||
transformer = None
|
||||
vae = None
|
||||
|
||||
if args.fp16 and args.bf16:
|
||||
raise ValueError("You cannot pass both --fp16 and --bf16 at the same time.")
|
||||
|
||||
dtype = torch.float16 if args.fp16 else torch.bfloat16 if args.bf16 else torch.float32
|
||||
|
||||
if args.transformer_ckpt_path is not None:
|
||||
transformer = convert_transformer(args.transformer_ckpt_path)
|
||||
transformer = convert_transformer(
|
||||
args.transformer_ckpt_path,
|
||||
args.num_layers,
|
||||
args.num_attention_heads,
|
||||
args.use_rotary_positional_embeddings,
|
||||
dtype,
|
||||
)
|
||||
if args.vae_ckpt_path is not None:
|
||||
vae = convert_vae(args.vae_ckpt_path)
|
||||
vae = convert_vae(args.vae_ckpt_path, args.scaling_factor, dtype)
|
||||
|
||||
text_encoder_id = "google/t5-v1_1-xxl"
|
||||
tokenizer = T5Tokenizer.from_pretrained(text_encoder_id, model_max_length=TOKENIZER_MAX_LENGTH)
|
||||
text_encoder = T5EncoderModel.from_pretrained(text_encoder_id, cache_dir=args.text_encoder_cache_dir)
|
||||
|
||||
# Apparently, the conversion does not work any more without this :shrug:
|
||||
for param in text_encoder.parameters():
|
||||
param.data = param.data.contiguous()
|
||||
|
||||
scheduler = CogVideoXDDIMScheduler.from_config(
|
||||
{
|
||||
"snr_shift_scale": 3.0,
|
||||
"snr_shift_scale": args.snr_shift_scale,
|
||||
"beta_end": 0.012,
|
||||
"beta_schedule": "scaled_linear",
|
||||
"beta_start": 0.00085,
|
||||
@@ -208,7 +249,7 @@ if __name__ == "__main__":
|
||||
"prediction_type": "v_prediction",
|
||||
"rescale_betas_zero_snr": True,
|
||||
"set_alpha_to_one": True,
|
||||
"timestep_spacing": "linspace",
|
||||
"timestep_spacing": "trailing",
|
||||
}
|
||||
)
|
||||
|
||||
@@ -218,5 +259,10 @@ if __name__ == "__main__":
|
||||
|
||||
if args.fp16:
|
||||
pipe = pipe.to(dtype=torch.float16)
|
||||
if args.bf16:
|
||||
pipe = pipe.to(dtype=torch.bfloat16)
|
||||
|
||||
# We don't use variant here because the model must be run in fp16 (2B) or bf16 (5B). It would be weird
|
||||
# for users to specify variant when the default is not fp32 and they want to run with the correct default (which
|
||||
# is either fp16/bf16 here).
|
||||
pipe.save_pretrained(args.output_path, safe_serialization=True, push_to_hub=args.push_to_hub)
|
||||
|
||||
@@ -254,7 +254,7 @@ version_range_max = max(sys.version_info[1], 10) + 1
|
||||
|
||||
setup(
|
||||
name="diffusers",
|
||||
version="0.30.0.dev0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
|
||||
version="0.30.2", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
|
||||
description="State-of-the-art diffusion in PyTorch and JAX.",
|
||||
long_description=open("README.md", "r", encoding="utf-8").read(),
|
||||
long_description_content_type="text/markdown",
|
||||
|
||||
@@ -1,4 +1,4 @@
|
||||
__version__ = "0.30.0.dev0"
|
||||
__version__ = "0.30.2"
|
||||
|
||||
from typing import TYPE_CHECKING
|
||||
|
||||
|
||||
@@ -208,6 +208,8 @@ class IPAdapterMixin:
|
||||
pretrained_model_name_or_path_or_dict,
|
||||
subfolder=image_encoder_subfolder,
|
||||
low_cpu_mem_usage=low_cpu_mem_usage,
|
||||
cache_dir=cache_dir,
|
||||
local_files_only=local_files_only,
|
||||
).to(self.device, dtype=self.dtype)
|
||||
self.register_modules(image_encoder=image_encoder)
|
||||
else:
|
||||
|
||||
@@ -1489,10 +1489,10 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
|
||||
@classmethod
|
||||
@validate_hf_hub_args
|
||||
# Copied from diffusers.loaders.lora_pipeline.SD3LoraLoaderMixin.lora_state_dict
|
||||
def lora_state_dict(
|
||||
cls,
|
||||
pretrained_model_name_or_path_or_dict: Union[str, Dict[str, torch.Tensor]],
|
||||
return_alphas: bool = False,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
@@ -1577,7 +1577,26 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
allow_pickle=allow_pickle,
|
||||
)
|
||||
|
||||
return state_dict
|
||||
# For state dicts like
|
||||
# https://huggingface.co/TheLastBen/Jon_Snow_Flux_LoRA
|
||||
keys = list(state_dict.keys())
|
||||
network_alphas = {}
|
||||
for k in keys:
|
||||
if "alpha" in k:
|
||||
alpha_value = state_dict.get(k)
|
||||
if (torch.is_tensor(alpha_value) and torch.is_floating_point(alpha_value)) or isinstance(
|
||||
alpha_value, float
|
||||
):
|
||||
network_alphas[k] = state_dict.pop(k)
|
||||
else:
|
||||
raise ValueError(
|
||||
f"The alpha key ({k}) seems to be incorrect. If you think this error is unexpected, please open as issue."
|
||||
)
|
||||
|
||||
if return_alphas:
|
||||
return state_dict, network_alphas
|
||||
else:
|
||||
return state_dict
|
||||
|
||||
def load_lora_weights(
|
||||
self, pretrained_model_name_or_path_or_dict: Union[str, Dict[str, torch.Tensor]], adapter_name=None, **kwargs
|
||||
@@ -1611,7 +1630,9 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
pretrained_model_name_or_path_or_dict = pretrained_model_name_or_path_or_dict.copy()
|
||||
|
||||
# First, ensure that the checkpoint is a compatible one and can be successfully loaded.
|
||||
state_dict = self.lora_state_dict(pretrained_model_name_or_path_or_dict, **kwargs)
|
||||
state_dict, network_alphas = self.lora_state_dict(
|
||||
pretrained_model_name_or_path_or_dict, return_alphas=True, **kwargs
|
||||
)
|
||||
|
||||
is_correct_format = all("lora" in key or "dora_scale" in key for key in state_dict.keys())
|
||||
if not is_correct_format:
|
||||
@@ -1619,6 +1640,7 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
|
||||
self.load_lora_into_transformer(
|
||||
state_dict,
|
||||
network_alphas=network_alphas,
|
||||
transformer=getattr(self, self.transformer_name) if not hasattr(self, "transformer") else self.transformer,
|
||||
adapter_name=adapter_name,
|
||||
_pipeline=self,
|
||||
@@ -1628,7 +1650,7 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
if len(text_encoder_state_dict) > 0:
|
||||
self.load_lora_into_text_encoder(
|
||||
text_encoder_state_dict,
|
||||
network_alphas=None,
|
||||
network_alphas=network_alphas,
|
||||
text_encoder=self.text_encoder,
|
||||
prefix="text_encoder",
|
||||
lora_scale=self.lora_scale,
|
||||
@@ -1637,8 +1659,7 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
)
|
||||
|
||||
@classmethod
|
||||
# Copied from diffusers.loaders.lora_pipeline.SD3LoraLoaderMixin.load_lora_into_transformer
|
||||
def load_lora_into_transformer(cls, state_dict, transformer, adapter_name=None, _pipeline=None):
|
||||
def load_lora_into_transformer(cls, state_dict, network_alphas, transformer, adapter_name=None, _pipeline=None):
|
||||
"""
|
||||
This will load the LoRA layers specified in `state_dict` into `transformer`.
|
||||
|
||||
@@ -1647,6 +1668,10 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
A standard state dict containing the lora layer parameters. The keys can either be indexed directly
|
||||
into the unet or prefixed with an additional `unet` which can be used to distinguish between text
|
||||
encoder lora layers.
|
||||
network_alphas (`Dict[str, float]`):
|
||||
The value of the network alpha used for stable learning and preventing underflow. This value has the
|
||||
same meaning as the `--network_alpha` option in the kohya-ss trainer script. Refer to [this
|
||||
link](https://github.com/darkstorm2150/sd-scripts/blob/main/docs/train_network_README-en.md#execute-learning).
|
||||
transformer (`SD3Transformer2DModel`):
|
||||
The Transformer model to load the LoRA layers into.
|
||||
adapter_name (`str`, *optional*):
|
||||
@@ -1678,7 +1703,12 @@ class FluxLoraLoaderMixin(LoraBaseMixin):
|
||||
if "lora_B" in key:
|
||||
rank[key] = val.shape[1]
|
||||
|
||||
lora_config_kwargs = get_peft_kwargs(rank, network_alpha_dict=None, peft_state_dict=state_dict)
|
||||
if network_alphas is not None and len(network_alphas) >= 1:
|
||||
prefix = cls.transformer_name
|
||||
alpha_keys = [k for k in network_alphas.keys() if k.startswith(prefix) and k.split(".")[0] == prefix]
|
||||
network_alphas = {k.replace(f"{prefix}.", ""): v for k, v in network_alphas.items() if k in alpha_keys}
|
||||
|
||||
lora_config_kwargs = get_peft_kwargs(rank, network_alpha_dict=network_alphas, peft_state_dict=state_dict)
|
||||
if "use_dora" in lora_config_kwargs:
|
||||
if lora_config_kwargs["use_dora"] and is_peft_version("<", "0.9.0"):
|
||||
raise ValueError(
|
||||
|
||||
@@ -23,6 +23,7 @@ from packaging import version
|
||||
from ..utils import deprecate, is_transformers_available, logging
|
||||
from .single_file_utils import (
|
||||
SingleFileComponentError,
|
||||
_is_legacy_scheduler_kwargs,
|
||||
_is_model_weights_in_cached_folder,
|
||||
_legacy_load_clip_tokenizer,
|
||||
_legacy_load_safety_checker,
|
||||
@@ -42,7 +43,6 @@ logger = logging.get_logger(__name__)
|
||||
# Legacy behaviour. `from_single_file` does not load the safety checker unless explicitly provided
|
||||
SINGLE_FILE_OPTIONAL_COMPONENTS = ["safety_checker"]
|
||||
|
||||
|
||||
if is_transformers_available():
|
||||
import transformers
|
||||
from transformers import PreTrainedModel, PreTrainedTokenizer
|
||||
@@ -135,7 +135,7 @@ def load_single_file_sub_model(
|
||||
class_obj, checkpoint=checkpoint, config=cached_model_config_path, local_files_only=local_files_only
|
||||
)
|
||||
|
||||
elif is_diffusers_scheduler and is_legacy_loading:
|
||||
elif is_diffusers_scheduler and (is_legacy_loading or _is_legacy_scheduler_kwargs(kwargs)):
|
||||
loaded_sub_model = _legacy_load_scheduler(
|
||||
class_obj, checkpoint=checkpoint, component_name=name, original_config=original_config, **kwargs
|
||||
)
|
||||
|
||||
@@ -79,7 +79,10 @@ CHECKPOINT_KEY_NAMES = {
|
||||
"animatediff_sdxl_beta": "up_blocks.2.motion_modules.0.temporal_transformer.norm.weight",
|
||||
"animatediff_scribble": "controlnet_cond_embedding.conv_in.weight",
|
||||
"animatediff_rgb": "controlnet_cond_embedding.weight",
|
||||
"flux": "double_blocks.0.img_attn.norm.key_norm.scale",
|
||||
"flux": [
|
||||
"double_blocks.0.img_attn.norm.key_norm.scale",
|
||||
"model.diffusion_model.double_blocks.0.img_attn.norm.key_norm.scale",
|
||||
],
|
||||
}
|
||||
|
||||
DIFFUSERS_DEFAULT_PIPELINE_PATHS = {
|
||||
@@ -88,11 +91,11 @@ DIFFUSERS_DEFAULT_PIPELINE_PATHS = {
|
||||
"xl_inpaint": {"pretrained_model_name_or_path": "diffusers/stable-diffusion-xl-1.0-inpainting-0.1"},
|
||||
"playground-v2-5": {"pretrained_model_name_or_path": "playgroundai/playground-v2.5-1024px-aesthetic"},
|
||||
"upscale": {"pretrained_model_name_or_path": "stabilityai/stable-diffusion-x4-upscaler"},
|
||||
"inpainting": {"pretrained_model_name_or_path": "runwayml/stable-diffusion-inpainting"},
|
||||
"inpainting": {"pretrained_model_name_or_path": "Lykon/dreamshaper-8-inpainting"},
|
||||
"inpainting_v2": {"pretrained_model_name_or_path": "stabilityai/stable-diffusion-2-inpainting"},
|
||||
"controlnet": {"pretrained_model_name_or_path": "lllyasviel/control_v11p_sd15_canny"},
|
||||
"v2": {"pretrained_model_name_or_path": "stabilityai/stable-diffusion-2-1"},
|
||||
"v1": {"pretrained_model_name_or_path": "runwayml/stable-diffusion-v1-5"},
|
||||
"v1": {"pretrained_model_name_or_path": "Lykon/dreamshaper-8"},
|
||||
"stable_cascade_stage_b": {"pretrained_model_name_or_path": "stabilityai/stable-cascade", "subfolder": "decoder"},
|
||||
"stable_cascade_stage_b_lite": {
|
||||
"pretrained_model_name_or_path": "stabilityai/stable-cascade",
|
||||
@@ -258,7 +261,7 @@ SCHEDULER_DEFAULT_CONFIG = {
|
||||
"timestep_spacing": "leading",
|
||||
}
|
||||
|
||||
LDM_VAE_KEY = "first_stage_model."
|
||||
LDM_VAE_KEYS = ["first_stage_model.", "vae."]
|
||||
LDM_VAE_DEFAULT_SCALING_FACTOR = 0.18215
|
||||
PLAYGROUND_VAE_SCALING_FACTOR = 0.5
|
||||
LDM_UNET_KEY = "model.diffusion_model."
|
||||
@@ -267,8 +270,8 @@ LDM_CLIP_PREFIX_TO_REMOVE = [
|
||||
"cond_stage_model.transformer.",
|
||||
"conditioner.embedders.0.transformer.",
|
||||
]
|
||||
OPEN_CLIP_PREFIX = "conditioner.embedders.0.model."
|
||||
LDM_OPEN_CLIP_TEXT_PROJECTION_DIM = 1024
|
||||
SCHEDULER_LEGACY_KWARGS = ["prediction_type", "scheduler_type"]
|
||||
|
||||
VALID_URL_PREFIXES = ["https://huggingface.co/", "huggingface.co/", "hf.co/", "https://hf.co/"]
|
||||
|
||||
@@ -318,6 +321,10 @@ def _is_model_weights_in_cached_folder(cached_folder, name):
|
||||
return weights_exist
|
||||
|
||||
|
||||
def _is_legacy_scheduler_kwargs(kwargs):
|
||||
return any(k in SCHEDULER_LEGACY_KWARGS for k in kwargs.keys())
|
||||
|
||||
|
||||
def load_single_file_checkpoint(
|
||||
pretrained_model_link_or_path,
|
||||
force_download=False,
|
||||
@@ -516,8 +523,10 @@ def infer_diffusers_model_type(checkpoint):
|
||||
else:
|
||||
model_type = "animatediff_v3"
|
||||
|
||||
elif CHECKPOINT_KEY_NAMES["flux"] in checkpoint:
|
||||
if "guidance_in.in_layer.bias" in checkpoint:
|
||||
elif any(key in checkpoint for key in CHECKPOINT_KEY_NAMES["flux"]):
|
||||
if any(
|
||||
g in checkpoint for g in ["guidance_in.in_layer.bias", "model.diffusion_model.guidance_in.in_layer.bias"]
|
||||
):
|
||||
model_type = "flux-dev"
|
||||
else:
|
||||
model_type = "flux-schnell"
|
||||
@@ -1176,7 +1185,11 @@ def convert_ldm_vae_checkpoint(checkpoint, config):
|
||||
# remove the LDM_VAE_KEY prefix from the ldm checkpoint keys so that it is easier to map them to diffusers keys
|
||||
vae_state_dict = {}
|
||||
keys = list(checkpoint.keys())
|
||||
vae_key = LDM_VAE_KEY if any(k.startswith(LDM_VAE_KEY) for k in keys) else ""
|
||||
vae_key = ""
|
||||
for ldm_vae_key in LDM_VAE_KEYS:
|
||||
if any(k.startswith(ldm_vae_key) for k in keys):
|
||||
vae_key = ldm_vae_key
|
||||
|
||||
for key in keys:
|
||||
if key.startswith(vae_key):
|
||||
vae_state_dict[key.replace(vae_key, "")] = checkpoint.get(key)
|
||||
@@ -1477,14 +1490,22 @@ def _legacy_load_scheduler(
|
||||
|
||||
if scheduler_type is not None:
|
||||
deprecation_message = (
|
||||
"Please pass an instance of a Scheduler object directly to the `scheduler` argument in `from_single_file`."
|
||||
"Please pass an instance of a Scheduler object directly to the `scheduler` argument in `from_single_file`\n\n"
|
||||
"Example:\n\n"
|
||||
"from diffusers import StableDiffusionPipeline, DDIMScheduler\n\n"
|
||||
"scheduler = DDIMScheduler()\n"
|
||||
"pipe = StableDiffusionPipeline.from_single_file(<checkpoint path>, scheduler=scheduler)\n"
|
||||
)
|
||||
deprecate("scheduler_type", "1.0.0", deprecation_message)
|
||||
|
||||
if prediction_type is not None:
|
||||
deprecation_message = (
|
||||
"Please configure an instance of a Scheduler with the appropriate `prediction_type` "
|
||||
"and pass the object directly to the `scheduler` argument in `from_single_file`."
|
||||
"Please configure an instance of a Scheduler with the appropriate `prediction_type` and "
|
||||
"pass the object directly to the `scheduler` argument in `from_single_file`.\n\n"
|
||||
"Example:\n\n"
|
||||
"from diffusers import StableDiffusionPipeline, DDIMScheduler\n\n"
|
||||
'scheduler = DDIMScheduler(prediction_type="v_prediction")\n'
|
||||
"pipe = StableDiffusionPipeline.from_single_file(<checkpoint path>, scheduler=scheduler)\n"
|
||||
)
|
||||
deprecate("prediction_type", "1.0.0", deprecation_message)
|
||||
|
||||
@@ -1881,6 +1902,10 @@ def convert_animatediff_checkpoint_to_diffusers(checkpoint, **kwargs):
|
||||
|
||||
def convert_flux_transformer_checkpoint_to_diffusers(checkpoint, **kwargs):
|
||||
converted_state_dict = {}
|
||||
keys = list(checkpoint.keys())
|
||||
for k in keys:
|
||||
if "model.diffusion_model." in k:
|
||||
checkpoint[k.replace("model.diffusion_model.", "")] = checkpoint.pop(k)
|
||||
|
||||
num_layers = list(set(int(k.split(".", 2)[1]) for k in checkpoint if "double_blocks." in k))[-1] + 1 # noqa: C401
|
||||
num_single_layers = list(set(int(k.split(".", 2)[1]) for k in checkpoint if "single_blocks." in k))[-1] + 1 # noqa: C401
|
||||
|
||||
@@ -1868,6 +1868,148 @@ class FluxAttnProcessor2_0:
|
||||
return hidden_states, encoder_hidden_states
|
||||
|
||||
|
||||
class CogVideoXAttnProcessor2_0:
|
||||
r"""
|
||||
Processor for implementing scaled dot-product attention for the CogVideoX model. It applies a rotary embedding on
|
||||
query and key vectors, but does not include spatial normalization.
|
||||
"""
|
||||
|
||||
def __init__(self):
|
||||
if not hasattr(F, "scaled_dot_product_attention"):
|
||||
raise ImportError("CogVideoXAttnProcessor requires PyTorch 2.0, to use it, please upgrade PyTorch to 2.0.")
|
||||
|
||||
def __call__(
|
||||
self,
|
||||
attn: Attention,
|
||||
hidden_states: torch.Tensor,
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
attention_mask: Optional[torch.Tensor] = None,
|
||||
image_rotary_emb: Optional[torch.Tensor] = None,
|
||||
) -> torch.Tensor:
|
||||
text_seq_length = encoder_hidden_states.size(1)
|
||||
|
||||
hidden_states = torch.cat([encoder_hidden_states, hidden_states], dim=1)
|
||||
|
||||
batch_size, sequence_length, _ = (
|
||||
hidden_states.shape if encoder_hidden_states is None else encoder_hidden_states.shape
|
||||
)
|
||||
|
||||
if attention_mask is not None:
|
||||
attention_mask = attn.prepare_attention_mask(attention_mask, sequence_length, batch_size)
|
||||
attention_mask = attention_mask.view(batch_size, attn.heads, -1, attention_mask.shape[-1])
|
||||
|
||||
query = attn.to_q(hidden_states)
|
||||
key = attn.to_k(hidden_states)
|
||||
value = attn.to_v(hidden_states)
|
||||
|
||||
inner_dim = key.shape[-1]
|
||||
head_dim = inner_dim // attn.heads
|
||||
|
||||
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
key = key.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
value = value.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
|
||||
if attn.norm_q is not None:
|
||||
query = attn.norm_q(query)
|
||||
if attn.norm_k is not None:
|
||||
key = attn.norm_k(key)
|
||||
|
||||
# Apply RoPE if needed
|
||||
if image_rotary_emb is not None:
|
||||
from .embeddings import apply_rotary_emb
|
||||
|
||||
query[:, :, text_seq_length:] = apply_rotary_emb(query[:, :, text_seq_length:], image_rotary_emb)
|
||||
if not attn.is_cross_attention:
|
||||
key[:, :, text_seq_length:] = apply_rotary_emb(key[:, :, text_seq_length:], image_rotary_emb)
|
||||
|
||||
hidden_states = F.scaled_dot_product_attention(
|
||||
query, key, value, attn_mask=attention_mask, dropout_p=0.0, is_causal=False
|
||||
)
|
||||
|
||||
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
|
||||
|
||||
# linear proj
|
||||
hidden_states = attn.to_out[0](hidden_states)
|
||||
# dropout
|
||||
hidden_states = attn.to_out[1](hidden_states)
|
||||
|
||||
encoder_hidden_states, hidden_states = hidden_states.split(
|
||||
[text_seq_length, hidden_states.size(1) - text_seq_length], dim=1
|
||||
)
|
||||
return hidden_states, encoder_hidden_states
|
||||
|
||||
|
||||
class FusedCogVideoXAttnProcessor2_0:
|
||||
r"""
|
||||
Processor for implementing scaled dot-product attention for the CogVideoX model. It applies a rotary embedding on
|
||||
query and key vectors, but does not include spatial normalization.
|
||||
"""
|
||||
|
||||
def __init__(self):
|
||||
if not hasattr(F, "scaled_dot_product_attention"):
|
||||
raise ImportError("CogVideoXAttnProcessor requires PyTorch 2.0, to use it, please upgrade PyTorch to 2.0.")
|
||||
|
||||
def __call__(
|
||||
self,
|
||||
attn: Attention,
|
||||
hidden_states: torch.Tensor,
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
attention_mask: Optional[torch.Tensor] = None,
|
||||
image_rotary_emb: Optional[torch.Tensor] = None,
|
||||
) -> torch.Tensor:
|
||||
text_seq_length = encoder_hidden_states.size(1)
|
||||
|
||||
hidden_states = torch.cat([encoder_hidden_states, hidden_states], dim=1)
|
||||
|
||||
batch_size, sequence_length, _ = (
|
||||
hidden_states.shape if encoder_hidden_states is None else encoder_hidden_states.shape
|
||||
)
|
||||
|
||||
if attention_mask is not None:
|
||||
attention_mask = attn.prepare_attention_mask(attention_mask, sequence_length, batch_size)
|
||||
attention_mask = attention_mask.view(batch_size, attn.heads, -1, attention_mask.shape[-1])
|
||||
|
||||
qkv = attn.to_qkv(hidden_states)
|
||||
split_size = qkv.shape[-1] // 3
|
||||
query, key, value = torch.split(qkv, split_size, dim=-1)
|
||||
|
||||
inner_dim = key.shape[-1]
|
||||
head_dim = inner_dim // attn.heads
|
||||
|
||||
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
key = key.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
value = value.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
|
||||
|
||||
if attn.norm_q is not None:
|
||||
query = attn.norm_q(query)
|
||||
if attn.norm_k is not None:
|
||||
key = attn.norm_k(key)
|
||||
|
||||
# Apply RoPE if needed
|
||||
if image_rotary_emb is not None:
|
||||
from .embeddings import apply_rotary_emb
|
||||
|
||||
query[:, :, text_seq_length:] = apply_rotary_emb(query[:, :, text_seq_length:], image_rotary_emb)
|
||||
if not attn.is_cross_attention:
|
||||
key[:, :, text_seq_length:] = apply_rotary_emb(key[:, :, text_seq_length:], image_rotary_emb)
|
||||
|
||||
hidden_states = F.scaled_dot_product_attention(
|
||||
query, key, value, attn_mask=attention_mask, dropout_p=0.0, is_causal=False
|
||||
)
|
||||
|
||||
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
|
||||
|
||||
# linear proj
|
||||
hidden_states = attn.to_out[0](hidden_states)
|
||||
# dropout
|
||||
hidden_states = attn.to_out[1](hidden_states)
|
||||
|
||||
encoder_hidden_states, hidden_states = hidden_states.split(
|
||||
[text_seq_length, hidden_states.size(1) - text_seq_length], dim=1
|
||||
)
|
||||
return hidden_states, encoder_hidden_states
|
||||
|
||||
|
||||
class XFormersAttnAddedKVProcessor:
|
||||
r"""
|
||||
Processor for implementing memory efficient attention using xFormers.
|
||||
|
||||
@@ -36,7 +36,7 @@ logger = logging.get_logger(__name__) # pylint: disable=invalid-name
|
||||
|
||||
|
||||
class CogVideoXSafeConv3d(nn.Conv3d):
|
||||
"""
|
||||
r"""
|
||||
A 3D convolution layer that splits the input tensor into smaller parts to avoid OOM in CogVideoX Model.
|
||||
"""
|
||||
|
||||
@@ -68,12 +68,12 @@ class CogVideoXCausalConv3d(nn.Module):
|
||||
r"""A 3D causal convolution layer that pads the input tensor to ensure causality in CogVideoX Model.
|
||||
|
||||
Args:
|
||||
in_channels (int): Number of channels in the input tensor.
|
||||
out_channels (int): Number of output channels.
|
||||
kernel_size (Union[int, Tuple[int, int, int]]): Size of the convolutional kernel.
|
||||
stride (int, optional): Stride of the convolution. Default is 1.
|
||||
dilation (int, optional): Dilation rate of the convolution. Default is 1.
|
||||
pad_mode (str, optional): Padding mode. Default is "constant".
|
||||
in_channels (`int`): Number of channels in the input tensor.
|
||||
out_channels (`int`): Number of output channels produced by the convolution.
|
||||
kernel_size (`int` or `Tuple[int, int, int]`): Kernel size of the convolutional kernel.
|
||||
stride (`int`, defaults to `1`): Stride of the convolution.
|
||||
dilation (`int`, defaults to `1`): Dilation rate of the convolution.
|
||||
pad_mode (`str`, defaults to `"constant"`): Padding mode.
|
||||
"""
|
||||
|
||||
def __init__(
|
||||
@@ -118,19 +118,12 @@ class CogVideoXCausalConv3d(nn.Module):
|
||||
self.conv_cache = None
|
||||
|
||||
def fake_context_parallel_forward(self, inputs: torch.Tensor) -> torch.Tensor:
|
||||
dim = self.temporal_dim
|
||||
kernel_size = self.time_kernel_size
|
||||
if kernel_size == 1:
|
||||
return inputs
|
||||
|
||||
inputs = inputs.transpose(0, dim)
|
||||
|
||||
if self.conv_cache is not None:
|
||||
inputs = torch.cat([self.conv_cache.transpose(0, dim).to(inputs.device), inputs], dim=0)
|
||||
else:
|
||||
inputs = torch.cat([inputs[:1]] * (kernel_size - 1) + [inputs], dim=0)
|
||||
|
||||
inputs = inputs.transpose(0, dim).contiguous()
|
||||
if kernel_size > 1:
|
||||
cached_inputs = (
|
||||
[self.conv_cache] if self.conv_cache is not None else [inputs[:, :, :1]] * (kernel_size - 1)
|
||||
)
|
||||
inputs = torch.cat(cached_inputs + [inputs], dim=2)
|
||||
return inputs
|
||||
|
||||
def _clear_fake_context_parallel_cache(self):
|
||||
@@ -138,16 +131,17 @@ class CogVideoXCausalConv3d(nn.Module):
|
||||
self.conv_cache = None
|
||||
|
||||
def forward(self, inputs: torch.Tensor) -> torch.Tensor:
|
||||
input_parallel = self.fake_context_parallel_forward(inputs)
|
||||
inputs = self.fake_context_parallel_forward(inputs)
|
||||
|
||||
self._clear_fake_context_parallel_cache()
|
||||
self.conv_cache = input_parallel[:, :, -self.time_kernel_size + 1 :].contiguous().detach().clone().cpu()
|
||||
# Note: we could move these to the cpu for a lower maximum memory usage but its only a few
|
||||
# hundred megabytes and so let's not do it for now
|
||||
self.conv_cache = inputs[:, :, -self.time_kernel_size + 1 :].clone()
|
||||
|
||||
padding_2d = (self.width_pad, self.width_pad, self.height_pad, self.height_pad)
|
||||
input_parallel = F.pad(input_parallel, padding_2d, mode="constant", value=0)
|
||||
inputs = F.pad(inputs, padding_2d, mode="constant", value=0)
|
||||
|
||||
output_parallel = self.conv(input_parallel)
|
||||
output = output_parallel
|
||||
output = self.conv(inputs)
|
||||
return output
|
||||
|
||||
|
||||
@@ -163,6 +157,8 @@ class CogVideoXSpatialNorm3D(nn.Module):
|
||||
The number of channels for input to group normalization layer, and output of the spatial norm layer.
|
||||
zq_channels (`int`):
|
||||
The number of channels for the quantized vector as described in the paper.
|
||||
groups (`int`):
|
||||
Number of groups to separate the channels into for group normalization.
|
||||
"""
|
||||
|
||||
def __init__(
|
||||
@@ -197,17 +193,26 @@ class CogVideoXResnetBlock3D(nn.Module):
|
||||
A 3D ResNet block used in the CogVideoX model.
|
||||
|
||||
Args:
|
||||
in_channels (int): Number of input channels.
|
||||
out_channels (Optional[int], optional):
|
||||
Number of output channels. If None, defaults to `in_channels`. Default is None.
|
||||
dropout (float, optional): Dropout rate. Default is 0.0.
|
||||
temb_channels (int, optional): Number of time embedding channels. Default is 512.
|
||||
groups (int, optional): Number of groups for group normalization. Default is 32.
|
||||
eps (float, optional): Epsilon value for normalization layers. Default is 1e-6.
|
||||
non_linearity (str, optional): Activation function to use. Default is "swish".
|
||||
conv_shortcut (bool, optional): If True, use a convolutional shortcut. Default is False.
|
||||
spatial_norm_dim (Optional[int], optional): Dimension of the spatial normalization. Default is None.
|
||||
pad_mode (str, optional): Padding mode. Default is "first".
|
||||
in_channels (`int`):
|
||||
Number of input channels.
|
||||
out_channels (`int`, *optional*):
|
||||
Number of output channels. If None, defaults to `in_channels`.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
Dropout rate.
|
||||
temb_channels (`int`, defaults to `512`):
|
||||
Number of time embedding channels.
|
||||
groups (`int`, defaults to `32`):
|
||||
Number of groups to separate the channels into for group normalization.
|
||||
eps (`float`, defaults to `1e-6`):
|
||||
Epsilon value for normalization layers.
|
||||
non_linearity (`str`, defaults to `"swish"`):
|
||||
Activation function to use.
|
||||
conv_shortcut (bool, defaults to `False`):
|
||||
Whether or not to use a convolution shortcut.
|
||||
spatial_norm_dim (`int`, *optional*):
|
||||
The dimension to use for spatial norm if it is to be used instead of group norm.
|
||||
pad_mode (str, defaults to `"first"`):
|
||||
Padding mode.
|
||||
"""
|
||||
|
||||
def __init__(
|
||||
@@ -309,18 +314,28 @@ class CogVideoXDownBlock3D(nn.Module):
|
||||
A downsampling block used in the CogVideoX model.
|
||||
|
||||
Args:
|
||||
in_channels (int): Number of input channels.
|
||||
out_channels (int): Number of output channels.
|
||||
temb_channels (int): Number of time embedding channels.
|
||||
dropout (float, optional): Dropout rate. Default is 0.0.
|
||||
num_layers (int, optional): Number of layers in the block. Default is 1.
|
||||
resnet_eps (float, optional): Epsilon value for the ResNet layers. Default is 1e-6.
|
||||
resnet_act_fn (str, optional): Activation function for the ResNet layers. Default is "swish".
|
||||
resnet_groups (int, optional): Number of groups for group normalization in the ResNet layers. Default is 32.
|
||||
add_downsample (bool, optional): If True, add a downsampling layer at the end of the block. Default is True.
|
||||
downsample_padding (int, optional): Padding for the downsampling layer. Default is 0.
|
||||
compress_time (bool, optional): If True, apply temporal compression. Default is False.
|
||||
pad_mode (str, optional): Padding mode. Default is "first".
|
||||
in_channels (`int`):
|
||||
Number of input channels.
|
||||
out_channels (`int`, *optional*):
|
||||
Number of output channels. If None, defaults to `in_channels`.
|
||||
temb_channels (`int`, defaults to `512`):
|
||||
Number of time embedding channels.
|
||||
num_layers (`int`, defaults to `1`):
|
||||
Number of resnet layers.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
Dropout rate.
|
||||
resnet_eps (`float`, defaults to `1e-6`):
|
||||
Epsilon value for normalization layers.
|
||||
resnet_act_fn (`str`, defaults to `"swish"`):
|
||||
Activation function to use.
|
||||
resnet_groups (`int`, defaults to `32`):
|
||||
Number of groups to separate the channels into for group normalization.
|
||||
add_downsample (`bool`, defaults to `True`):
|
||||
Whether or not to use a downsampling layer. If not used, output dimension would be same as input dimension.
|
||||
compress_time (`bool`, defaults to `False`):
|
||||
Whether or not to downsample across temporal dimension.
|
||||
pad_mode (str, defaults to `"first"`):
|
||||
Padding mode.
|
||||
"""
|
||||
|
||||
_supports_gradient_checkpointing = True
|
||||
@@ -405,15 +420,24 @@ class CogVideoXMidBlock3D(nn.Module):
|
||||
A middle block used in the CogVideoX model.
|
||||
|
||||
Args:
|
||||
in_channels (int): Number of input channels.
|
||||
temb_channels (int): Number of time embedding channels.
|
||||
dropout (float, optional): Dropout rate. Default is 0.0.
|
||||
num_layers (int, optional): Number of layers in the block. Default is 1.
|
||||
resnet_eps (float, optional): Epsilon value for the ResNet layers. Default is 1e-6.
|
||||
resnet_act_fn (str, optional): Activation function for the ResNet layers. Default is "swish".
|
||||
resnet_groups (int, optional): Number of groups for group normalization in the ResNet layers. Default is 32.
|
||||
spatial_norm_dim (Optional[int], optional): Dimension of the spatial normalization. Default is None.
|
||||
pad_mode (str, optional): Padding mode. Default is "first".
|
||||
in_channels (`int`):
|
||||
Number of input channels.
|
||||
temb_channels (`int`, defaults to `512`):
|
||||
Number of time embedding channels.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
Dropout rate.
|
||||
num_layers (`int`, defaults to `1`):
|
||||
Number of resnet layers.
|
||||
resnet_eps (`float`, defaults to `1e-6`):
|
||||
Epsilon value for normalization layers.
|
||||
resnet_act_fn (`str`, defaults to `"swish"`):
|
||||
Activation function to use.
|
||||
resnet_groups (`int`, defaults to `32`):
|
||||
Number of groups to separate the channels into for group normalization.
|
||||
spatial_norm_dim (`int`, *optional*):
|
||||
The dimension to use for spatial norm if it is to be used instead of group norm.
|
||||
pad_mode (str, defaults to `"first"`):
|
||||
Padding mode.
|
||||
"""
|
||||
|
||||
_supports_gradient_checkpointing = True
|
||||
@@ -480,19 +504,30 @@ class CogVideoXUpBlock3D(nn.Module):
|
||||
An upsampling block used in the CogVideoX model.
|
||||
|
||||
Args:
|
||||
in_channels (int): Number of input channels.
|
||||
out_channels (int): Number of output channels.
|
||||
temb_channels (int): Number of time embedding channels.
|
||||
dropout (float, optional): Dropout rate. Default is 0.0.
|
||||
num_layers (int, optional): Number of layers in the block. Default is 1.
|
||||
resnet_eps (float, optional): Epsilon value for the ResNet layers. Default is 1e-6.
|
||||
resnet_act_fn (str, optional): Activation function for the ResNet layers. Default is "swish".
|
||||
resnet_groups (int, optional): Number of groups for group normalization in the ResNet layers. Default is 32.
|
||||
spatial_norm_dim (int, optional): Dimension of the spatial normalization. Default is 16.
|
||||
add_upsample (bool, optional): If True, add an upsampling layer at the end of the block. Default is True.
|
||||
upsample_padding (int, optional): Padding for the upsampling layer. Default is 1.
|
||||
compress_time (bool, optional): If True, apply temporal compression. Default is False.
|
||||
pad_mode (str, optional): Padding mode. Default is "first".
|
||||
in_channels (`int`):
|
||||
Number of input channels.
|
||||
out_channels (`int`, *optional*):
|
||||
Number of output channels. If None, defaults to `in_channels`.
|
||||
temb_channels (`int`, defaults to `512`):
|
||||
Number of time embedding channels.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
Dropout rate.
|
||||
num_layers (`int`, defaults to `1`):
|
||||
Number of resnet layers.
|
||||
resnet_eps (`float`, defaults to `1e-6`):
|
||||
Epsilon value for normalization layers.
|
||||
resnet_act_fn (`str`, defaults to `"swish"`):
|
||||
Activation function to use.
|
||||
resnet_groups (`int`, defaults to `32`):
|
||||
Number of groups to separate the channels into for group normalization.
|
||||
spatial_norm_dim (`int`, defaults to `16`):
|
||||
The dimension to use for spatial norm if it is to be used instead of group norm.
|
||||
add_upsample (`bool`, defaults to `True`):
|
||||
Whether or not to use a upsampling layer. If not used, output dimension would be same as input dimension.
|
||||
compress_time (`bool`, defaults to `False`):
|
||||
Whether or not to downsample across temporal dimension.
|
||||
pad_mode (str, defaults to `"first"`):
|
||||
Padding mode.
|
||||
"""
|
||||
|
||||
def __init__(
|
||||
@@ -587,14 +622,12 @@ class CogVideoXEncoder3D(nn.Module):
|
||||
options.
|
||||
block_out_channels (`Tuple[int, ...]`, *optional*, defaults to `(64,)`):
|
||||
The number of output channels for each block.
|
||||
act_fn (`str`, *optional*, defaults to `"silu"`):
|
||||
The activation function to use. See `~diffusers.models.activations.get_activation` for available options.
|
||||
layers_per_block (`int`, *optional*, defaults to 2):
|
||||
The number of layers per block.
|
||||
norm_num_groups (`int`, *optional*, defaults to 32):
|
||||
The number of groups for normalization.
|
||||
act_fn (`str`, *optional*, defaults to `"silu"`):
|
||||
The activation function to use. See `~diffusers.models.activations.get_activation` for available options.
|
||||
double_z (`bool`, *optional*, defaults to `True`):
|
||||
Whether to double the number of output channels for the last block.
|
||||
"""
|
||||
|
||||
_supports_gradient_checkpointing = True
|
||||
@@ -723,14 +756,12 @@ class CogVideoXDecoder3D(nn.Module):
|
||||
The types of up blocks to use. See `~diffusers.models.unet_2d_blocks.get_up_block` for available options.
|
||||
block_out_channels (`Tuple[int, ...]`, *optional*, defaults to `(64,)`):
|
||||
The number of output channels for each block.
|
||||
act_fn (`str`, *optional*, defaults to `"silu"`):
|
||||
The activation function to use. See `~diffusers.models.activations.get_activation` for available options.
|
||||
layers_per_block (`int`, *optional*, defaults to 2):
|
||||
The number of layers per block.
|
||||
norm_num_groups (`int`, *optional*, defaults to 32):
|
||||
The number of groups for normalization.
|
||||
act_fn (`str`, *optional*, defaults to `"silu"`):
|
||||
The activation function to use. See `~diffusers.models.activations.get_activation` for available options.
|
||||
norm_type (`str`, *optional*, defaults to `"group"`):
|
||||
The normalization type to use. Can be either `"group"` or `"spatial"`.
|
||||
"""
|
||||
|
||||
_supports_gradient_checkpointing = True
|
||||
@@ -871,7 +902,7 @@ class AutoencoderKLCogVideoX(ModelMixin, ConfigMixin, FromOriginalModelMixin):
|
||||
Tuple of block output channels.
|
||||
act_fn (`str`, *optional*, defaults to `"silu"`): The activation function to use.
|
||||
sample_size (`int`, *optional*, defaults to `32`): Sample input size.
|
||||
scaling_factor (`float`, *optional*, defaults to 0.18215):
|
||||
scaling_factor (`float`, *optional*, defaults to `1.15258426`):
|
||||
The component-wise standard deviation of the trained latent space computed using the first batch of the
|
||||
training set. This is used to scale the latent space to have unit variance when training the diffusion
|
||||
model. The latents are scaled with the formula `z = z * scaling_factor` before being passed to the
|
||||
@@ -911,7 +942,8 @@ class AutoencoderKLCogVideoX(ModelMixin, ConfigMixin, FromOriginalModelMixin):
|
||||
norm_eps: float = 1e-6,
|
||||
norm_num_groups: int = 32,
|
||||
temporal_compression_ratio: float = 4,
|
||||
sample_size: int = 256,
|
||||
sample_height: int = 480,
|
||||
sample_width: int = 720,
|
||||
scaling_factor: float = 1.15258426,
|
||||
shift_factor: Optional[float] = None,
|
||||
latents_mean: Optional[Tuple[float]] = None,
|
||||
@@ -950,25 +982,105 @@ class AutoencoderKLCogVideoX(ModelMixin, ConfigMixin, FromOriginalModelMixin):
|
||||
self.use_slicing = False
|
||||
self.use_tiling = False
|
||||
|
||||
self.tile_sample_min_size = self.config.sample_size
|
||||
sample_size = (
|
||||
self.config.sample_size[0]
|
||||
if isinstance(self.config.sample_size, (list, tuple))
|
||||
else self.config.sample_size
|
||||
# Can be increased to decode more latent frames at once, but comes at a reasonable memory cost and it is not
|
||||
# recommended because the temporal parts of the VAE, here, are tricky to understand.
|
||||
# If you decode X latent frames together, the number of output frames is:
|
||||
# (X + (2 conv cache) + (2 time upscale_1) + (4 time upscale_2) - (2 causal conv downscale)) => X + 6 frames
|
||||
#
|
||||
# Example with num_latent_frames_batch_size = 2:
|
||||
# - 12 latent frames: (0, 1), (2, 3), (4, 5), (6, 7), (8, 9), (10, 11) are processed together
|
||||
# => (12 // 2 frame slices) * ((2 num_latent_frames_batch_size) + (2 conv cache) + (2 time upscale_1) + (4 time upscale_2) - (2 causal conv downscale))
|
||||
# => 6 * 8 = 48 frames
|
||||
# - 13 latent frames: (0, 1, 2) (special case), (3, 4), (5, 6), (7, 8), (9, 10), (11, 12) are processed together
|
||||
# => (1 frame slice) * ((3 num_latent_frames_batch_size) + (2 conv cache) + (2 time upscale_1) + (4 time upscale_2) - (2 causal conv downscale)) +
|
||||
# ((13 - 3) // 2) * ((2 num_latent_frames_batch_size) + (2 conv cache) + (2 time upscale_1) + (4 time upscale_2) - (2 causal conv downscale))
|
||||
# => 1 * 9 + 5 * 8 = 49 frames
|
||||
# It has been implemented this way so as to not have "magic values" in the code base that would be hard to explain. Note that
|
||||
# setting it to anything other than 2 would give poor results because the VAE hasn't been trained to be adaptive with different
|
||||
# number of temporal frames.
|
||||
self.num_latent_frames_batch_size = 2
|
||||
|
||||
# We make the minimum height and width of sample for tiling half that of the generally supported
|
||||
self.tile_sample_min_height = sample_height // 2
|
||||
self.tile_sample_min_width = sample_width // 2
|
||||
self.tile_latent_min_height = int(
|
||||
self.tile_sample_min_height / (2 ** (len(self.config.block_out_channels) - 1))
|
||||
)
|
||||
self.tile_latent_min_size = int(sample_size / (2 ** (len(self.config.block_out_channels) - 1)))
|
||||
self.tile_overlap_factor = 0.25
|
||||
self.tile_latent_min_width = int(self.tile_sample_min_width / (2 ** (len(self.config.block_out_channels) - 1)))
|
||||
|
||||
# These are experimental overlap factors that were chosen based on experimentation and seem to work best for
|
||||
# 720x480 (WxH) resolution. The above resolution is the strongly recommended generation resolution in CogVideoX
|
||||
# and so the tiling implementation has only been tested on those specific resolutions.
|
||||
self.tile_overlap_factor_height = 1 / 6
|
||||
self.tile_overlap_factor_width = 1 / 5
|
||||
|
||||
def _set_gradient_checkpointing(self, module, value=False):
|
||||
if isinstance(module, (CogVideoXEncoder3D, CogVideoXDecoder3D)):
|
||||
module.gradient_checkpointing = value
|
||||
|
||||
def clear_fake_context_parallel_cache(self):
|
||||
def _clear_fake_context_parallel_cache(self):
|
||||
for name, module in self.named_modules():
|
||||
if isinstance(module, CogVideoXCausalConv3d):
|
||||
logger.debug(f"Clearing fake Context Parallel cache for layer: {name}")
|
||||
module._clear_fake_context_parallel_cache()
|
||||
|
||||
def enable_tiling(
|
||||
self,
|
||||
tile_sample_min_height: Optional[int] = None,
|
||||
tile_sample_min_width: Optional[int] = None,
|
||||
tile_overlap_factor_height: Optional[float] = None,
|
||||
tile_overlap_factor_width: Optional[float] = None,
|
||||
) -> None:
|
||||
r"""
|
||||
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to
|
||||
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow
|
||||
processing larger images.
|
||||
|
||||
Args:
|
||||
tile_sample_min_height (`int`, *optional*):
|
||||
The minimum height required for a sample to be separated into tiles across the height dimension.
|
||||
tile_sample_min_width (`int`, *optional*):
|
||||
The minimum width required for a sample to be separated into tiles across the width dimension.
|
||||
tile_overlap_factor_height (`int`, *optional*):
|
||||
The minimum amount of overlap between two consecutive vertical tiles. This is to ensure that there are
|
||||
no tiling artifacts produced across the height dimension. Must be between 0 and 1. Setting a higher
|
||||
value might cause more tiles to be processed leading to slow down of the decoding process.
|
||||
tile_overlap_factor_width (`int`, *optional*):
|
||||
The minimum amount of overlap between two consecutive horizontal tiles. This is to ensure that there
|
||||
are no tiling artifacts produced across the width dimension. Must be between 0 and 1. Setting a higher
|
||||
value might cause more tiles to be processed leading to slow down of the decoding process.
|
||||
"""
|
||||
self.use_tiling = True
|
||||
self.tile_sample_min_height = tile_sample_min_height or self.tile_sample_min_height
|
||||
self.tile_sample_min_width = tile_sample_min_width or self.tile_sample_min_width
|
||||
self.tile_latent_min_height = int(
|
||||
self.tile_sample_min_height / (2 ** (len(self.config.block_out_channels) - 1))
|
||||
)
|
||||
self.tile_latent_min_width = int(self.tile_sample_min_width / (2 ** (len(self.config.block_out_channels) - 1)))
|
||||
self.tile_overlap_factor_height = tile_overlap_factor_height or self.tile_overlap_factor_height
|
||||
self.tile_overlap_factor_width = tile_overlap_factor_width or self.tile_overlap_factor_width
|
||||
|
||||
def disable_tiling(self) -> None:
|
||||
r"""
|
||||
Disable tiled VAE decoding. If `enable_tiling` was previously enabled, this method will go back to computing
|
||||
decoding in one step.
|
||||
"""
|
||||
self.use_tiling = False
|
||||
|
||||
def enable_slicing(self) -> None:
|
||||
r"""
|
||||
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to
|
||||
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
|
||||
"""
|
||||
self.use_slicing = True
|
||||
|
||||
def disable_slicing(self) -> None:
|
||||
r"""
|
||||
Disable sliced VAE decoding. If `enable_slicing` was previously enabled, this method will go back to computing
|
||||
decoding in one step.
|
||||
"""
|
||||
self.use_slicing = False
|
||||
|
||||
@apply_forward_hook
|
||||
def encode(
|
||||
self, x: torch.Tensor, return_dict: bool = True
|
||||
@@ -993,8 +1105,34 @@ class AutoencoderKLCogVideoX(ModelMixin, ConfigMixin, FromOriginalModelMixin):
|
||||
return (posterior,)
|
||||
return AutoencoderKLOutput(latent_dist=posterior)
|
||||
|
||||
def _decode(self, z: torch.Tensor, return_dict: bool = True) -> Union[DecoderOutput, torch.Tensor]:
|
||||
batch_size, num_channels, num_frames, height, width = z.shape
|
||||
|
||||
if self.use_tiling and (width > self.tile_latent_min_width or height > self.tile_latent_min_height):
|
||||
return self.tiled_decode(z, return_dict=return_dict)
|
||||
|
||||
frame_batch_size = self.num_latent_frames_batch_size
|
||||
dec = []
|
||||
for i in range(num_frames // frame_batch_size):
|
||||
remaining_frames = num_frames % frame_batch_size
|
||||
start_frame = frame_batch_size * i + (0 if i == 0 else remaining_frames)
|
||||
end_frame = frame_batch_size * (i + 1) + remaining_frames
|
||||
z_intermediate = z[:, :, start_frame:end_frame]
|
||||
if self.post_quant_conv is not None:
|
||||
z_intermediate = self.post_quant_conv(z_intermediate)
|
||||
z_intermediate = self.decoder(z_intermediate)
|
||||
dec.append(z_intermediate)
|
||||
|
||||
self._clear_fake_context_parallel_cache()
|
||||
dec = torch.cat(dec, dim=2)
|
||||
|
||||
if not return_dict:
|
||||
return (dec,)
|
||||
|
||||
return DecoderOutput(sample=dec)
|
||||
|
||||
@apply_forward_hook
|
||||
def decode(self, z: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, torch.FloatTensor]:
|
||||
def decode(self, z: torch.Tensor, return_dict: bool = True) -> Union[DecoderOutput, torch.Tensor]:
|
||||
"""
|
||||
Decode a batch of images.
|
||||
|
||||
@@ -1007,13 +1145,111 @@ class AutoencoderKLCogVideoX(ModelMixin, ConfigMixin, FromOriginalModelMixin):
|
||||
[`~models.vae.DecoderOutput`] or `tuple`:
|
||||
If return_dict is True, a [`~models.vae.DecoderOutput`] is returned, otherwise a plain `tuple` is
|
||||
returned.
|
||||
|
||||
"""
|
||||
if self.post_quant_conv is not None:
|
||||
z = self.post_quant_conv(z)
|
||||
dec = self.decoder(z)
|
||||
if self.use_slicing and z.shape[0] > 1:
|
||||
decoded_slices = [self._decode(z_slice).sample for z_slice in z.split(1)]
|
||||
decoded = torch.cat(decoded_slices)
|
||||
else:
|
||||
decoded = self._decode(z).sample
|
||||
|
||||
if not return_dict:
|
||||
return (decoded,)
|
||||
return DecoderOutput(sample=decoded)
|
||||
|
||||
def blend_v(self, a: torch.Tensor, b: torch.Tensor, blend_extent: int) -> torch.Tensor:
|
||||
blend_extent = min(a.shape[3], b.shape[3], blend_extent)
|
||||
for y in range(blend_extent):
|
||||
b[:, :, :, y, :] = a[:, :, :, -blend_extent + y, :] * (1 - y / blend_extent) + b[:, :, :, y, :] * (
|
||||
y / blend_extent
|
||||
)
|
||||
return b
|
||||
|
||||
def blend_h(self, a: torch.Tensor, b: torch.Tensor, blend_extent: int) -> torch.Tensor:
|
||||
blend_extent = min(a.shape[4], b.shape[4], blend_extent)
|
||||
for x in range(blend_extent):
|
||||
b[:, :, :, :, x] = a[:, :, :, :, -blend_extent + x] * (1 - x / blend_extent) + b[:, :, :, :, x] * (
|
||||
x / blend_extent
|
||||
)
|
||||
return b
|
||||
|
||||
def tiled_decode(self, z: torch.Tensor, return_dict: bool = True) -> Union[DecoderOutput, torch.Tensor]:
|
||||
r"""
|
||||
Decode a batch of images using a tiled decoder.
|
||||
|
||||
Args:
|
||||
z (`torch.Tensor`): Input batch of latent vectors.
|
||||
return_dict (`bool`, *optional*, defaults to `True`):
|
||||
Whether or not to return a [`~models.vae.DecoderOutput`] instead of a plain tuple.
|
||||
|
||||
Returns:
|
||||
[`~models.vae.DecoderOutput`] or `tuple`:
|
||||
If return_dict is True, a [`~models.vae.DecoderOutput`] is returned, otherwise a plain `tuple` is
|
||||
returned.
|
||||
"""
|
||||
# Rough memory assessment:
|
||||
# - In CogVideoX-2B, there are a total of 24 CausalConv3d layers.
|
||||
# - The biggest intermediate dimensions are: [1, 128, 9, 480, 720].
|
||||
# - Assume fp16 (2 bytes per value).
|
||||
# Memory required: 1 * 128 * 9 * 480 * 720 * 24 * 2 / 1024**3 = 17.8 GB
|
||||
#
|
||||
# Memory assessment when using tiling:
|
||||
# - Assume everything as above but now HxW is 240x360 by tiling in half
|
||||
# Memory required: 1 * 128 * 9 * 240 * 360 * 24 * 2 / 1024**3 = 4.5 GB
|
||||
|
||||
batch_size, num_channels, num_frames, height, width = z.shape
|
||||
|
||||
overlap_height = int(self.tile_latent_min_height * (1 - self.tile_overlap_factor_height))
|
||||
overlap_width = int(self.tile_latent_min_width * (1 - self.tile_overlap_factor_width))
|
||||
blend_extent_height = int(self.tile_sample_min_height * self.tile_overlap_factor_height)
|
||||
blend_extent_width = int(self.tile_sample_min_width * self.tile_overlap_factor_width)
|
||||
row_limit_height = self.tile_sample_min_height - blend_extent_height
|
||||
row_limit_width = self.tile_sample_min_width - blend_extent_width
|
||||
frame_batch_size = self.num_latent_frames_batch_size
|
||||
|
||||
# Split z into overlapping tiles and decode them separately.
|
||||
# The tiles have an overlap to avoid seams between tiles.
|
||||
rows = []
|
||||
for i in range(0, height, overlap_height):
|
||||
row = []
|
||||
for j in range(0, width, overlap_width):
|
||||
time = []
|
||||
for k in range(num_frames // frame_batch_size):
|
||||
remaining_frames = num_frames % frame_batch_size
|
||||
start_frame = frame_batch_size * k + (0 if k == 0 else remaining_frames)
|
||||
end_frame = frame_batch_size * (k + 1) + remaining_frames
|
||||
tile = z[
|
||||
:,
|
||||
:,
|
||||
start_frame:end_frame,
|
||||
i : i + self.tile_latent_min_height,
|
||||
j : j + self.tile_latent_min_width,
|
||||
]
|
||||
if self.post_quant_conv is not None:
|
||||
tile = self.post_quant_conv(tile)
|
||||
tile = self.decoder(tile)
|
||||
time.append(tile)
|
||||
self._clear_fake_context_parallel_cache()
|
||||
row.append(torch.cat(time, dim=2))
|
||||
rows.append(row)
|
||||
|
||||
result_rows = []
|
||||
for i, row in enumerate(rows):
|
||||
result_row = []
|
||||
for j, tile in enumerate(row):
|
||||
# blend the above tile and the left tile
|
||||
# to the current tile and add the current tile to the result row
|
||||
if i > 0:
|
||||
tile = self.blend_v(rows[i - 1][j], tile, blend_extent_height)
|
||||
if j > 0:
|
||||
tile = self.blend_h(row[j - 1], tile, blend_extent_width)
|
||||
result_row.append(tile[:, :, :, :row_limit_height, :row_limit_width])
|
||||
result_rows.append(torch.cat(result_row, dim=4))
|
||||
|
||||
dec = torch.cat(result_rows, dim=3)
|
||||
|
||||
if not return_dict:
|
||||
return (dec,)
|
||||
|
||||
return DecoderOutput(sample=dec)
|
||||
|
||||
def forward(
|
||||
|
||||
@@ -374,6 +374,90 @@ class CogVideoXPatchEmbed(nn.Module):
|
||||
return embeds
|
||||
|
||||
|
||||
def get_3d_rotary_pos_embed(
|
||||
embed_dim, crops_coords, grid_size, temporal_size, theta: int = 10000, use_real: bool = True
|
||||
) -> Union[torch.Tensor, Tuple[torch.Tensor, torch.Tensor]]:
|
||||
"""
|
||||
RoPE for video tokens with 3D structure.
|
||||
|
||||
Args:
|
||||
embed_dim: (`int`):
|
||||
The embedding dimension size, corresponding to hidden_size_head.
|
||||
crops_coords (`Tuple[int]`):
|
||||
The top-left and bottom-right coordinates of the crop.
|
||||
grid_size (`Tuple[int]`):
|
||||
The grid size of the spatial positional embedding (height, width).
|
||||
temporal_size (`int`):
|
||||
The size of the temporal dimension.
|
||||
theta (`float`):
|
||||
Scaling factor for frequency computation.
|
||||
use_real (`bool`):
|
||||
If True, return real part and imaginary part separately. Otherwise, return complex numbers.
|
||||
|
||||
Returns:
|
||||
`torch.Tensor`: positional embedding with shape `(temporal_size * grid_size[0] * grid_size[1], embed_dim/2)`.
|
||||
"""
|
||||
start, stop = crops_coords
|
||||
grid_h = np.linspace(start[0], stop[0], grid_size[0], endpoint=False, dtype=np.float32)
|
||||
grid_w = np.linspace(start[1], stop[1], grid_size[1], endpoint=False, dtype=np.float32)
|
||||
grid_t = np.linspace(0, temporal_size, temporal_size, endpoint=False, dtype=np.float32)
|
||||
|
||||
# Compute dimensions for each axis
|
||||
dim_t = embed_dim // 4
|
||||
dim_h = embed_dim // 8 * 3
|
||||
dim_w = embed_dim // 8 * 3
|
||||
|
||||
# Temporal frequencies
|
||||
freqs_t = 1.0 / (theta ** (torch.arange(0, dim_t, 2).float() / dim_t))
|
||||
grid_t = torch.from_numpy(grid_t).float()
|
||||
freqs_t = torch.einsum("n , f -> n f", grid_t, freqs_t)
|
||||
freqs_t = freqs_t.repeat_interleave(2, dim=-1)
|
||||
|
||||
# Spatial frequencies for height and width
|
||||
freqs_h = 1.0 / (theta ** (torch.arange(0, dim_h, 2).float() / dim_h))
|
||||
freqs_w = 1.0 / (theta ** (torch.arange(0, dim_w, 2).float() / dim_w))
|
||||
grid_h = torch.from_numpy(grid_h).float()
|
||||
grid_w = torch.from_numpy(grid_w).float()
|
||||
freqs_h = torch.einsum("n , f -> n f", grid_h, freqs_h)
|
||||
freqs_w = torch.einsum("n , f -> n f", grid_w, freqs_w)
|
||||
freqs_h = freqs_h.repeat_interleave(2, dim=-1)
|
||||
freqs_w = freqs_w.repeat_interleave(2, dim=-1)
|
||||
|
||||
# Broadcast and concatenate tensors along specified dimension
|
||||
def broadcast(tensors, dim=-1):
|
||||
num_tensors = len(tensors)
|
||||
shape_lens = {len(t.shape) for t in tensors}
|
||||
assert len(shape_lens) == 1, "tensors must all have the same number of dimensions"
|
||||
shape_len = list(shape_lens)[0]
|
||||
dim = (dim + shape_len) if dim < 0 else dim
|
||||
dims = list(zip(*(list(t.shape) for t in tensors)))
|
||||
expandable_dims = [(i, val) for i, val in enumerate(dims) if i != dim]
|
||||
assert all(
|
||||
[*(len(set(t[1])) <= 2 for t in expandable_dims)]
|
||||
), "invalid dimensions for broadcastable concatenation"
|
||||
max_dims = [(t[0], max(t[1])) for t in expandable_dims]
|
||||
expanded_dims = [(t[0], (t[1],) * num_tensors) for t in max_dims]
|
||||
expanded_dims.insert(dim, (dim, dims[dim]))
|
||||
expandable_shapes = list(zip(*(t[1] for t in expanded_dims)))
|
||||
tensors = [t[0].expand(*t[1]) for t in zip(tensors, expandable_shapes)]
|
||||
return torch.cat(tensors, dim=dim)
|
||||
|
||||
freqs = broadcast((freqs_t[:, None, None, :], freqs_h[None, :, None, :], freqs_w[None, None, :, :]), dim=-1)
|
||||
|
||||
t, h, w, d = freqs.shape
|
||||
freqs = freqs.view(t * h * w, d)
|
||||
|
||||
# Generate sine and cosine components
|
||||
sin = freqs.sin()
|
||||
cos = freqs.cos()
|
||||
|
||||
if use_real:
|
||||
return cos, sin
|
||||
else:
|
||||
freqs_cis = torch.polar(torch.ones_like(freqs), freqs)
|
||||
return freqs_cis
|
||||
|
||||
|
||||
def get_2d_rotary_pos_embed(embed_dim, crops_coords, grid_size, use_real=True):
|
||||
"""
|
||||
RoPE for image tokens with 2d structure.
|
||||
|
||||
@@ -263,41 +263,6 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
|
||||
"""
|
||||
self.set_use_memory_efficient_attention_xformers(False)
|
||||
|
||||
def enable_dynamic_upcasting(self, upcast_dtype=None):
|
||||
upcast_dtype = upcast_dtype or torch.float32
|
||||
downcast_dtype = self.dtype
|
||||
|
||||
def upcast_hook_fn(module):
|
||||
module = module.to(upcast_dtype)
|
||||
|
||||
def downcast_hook_fn(module):
|
||||
module = module.to(downcast_dtype)
|
||||
|
||||
def fn_recursive_upcast(module):
|
||||
has_children = list(module.children())
|
||||
if not has_children:
|
||||
module.register_forward_pre_hook(upcast_hook_fn)
|
||||
module.register_forward_hook(downcast_hook_fn)
|
||||
|
||||
for child in module.children():
|
||||
fn_recursive_upcast(child)
|
||||
|
||||
for module in self.children():
|
||||
fn_recursive_upcast(module)
|
||||
|
||||
def disable_dynamic_upcasting(self):
|
||||
def fn_recursive_upcast(module):
|
||||
has_children = list(module.children())
|
||||
if not has_children:
|
||||
module._forward_pre_hooks = OrderedDict()
|
||||
module._forward_hooks = OrderedDict()
|
||||
|
||||
for child in module.children():
|
||||
fn_recursive_upcast(child)
|
||||
|
||||
for module in self.children():
|
||||
fn_recursive_upcast(module)
|
||||
|
||||
def save_pretrained(
|
||||
self,
|
||||
save_directory: Union[str, os.PathLike],
|
||||
|
||||
@@ -68,6 +68,21 @@ class AuraFlowPatchEmbed(nn.Module):
|
||||
self.height, self.width = height // patch_size, width // patch_size
|
||||
self.base_size = height // patch_size
|
||||
|
||||
def pe_selection_index_based_on_dim(self, h, w):
|
||||
# select subset of positional embedding based on H, W, where H, W is size of latent
|
||||
# PE will be viewed as 2d-grid, and H/p x W/p of the PE will be selected
|
||||
# because original input are in flattened format, we have to flatten this 2d grid as well.
|
||||
h_p, w_p = h // self.patch_size, w // self.patch_size
|
||||
original_pe_indexes = torch.arange(self.pos_embed.shape[1])
|
||||
h_max, w_max = int(self.pos_embed_max_size**0.5), int(self.pos_embed_max_size**0.5)
|
||||
original_pe_indexes = original_pe_indexes.view(h_max, w_max)
|
||||
starth = h_max // 2 - h_p // 2
|
||||
endh = starth + h_p
|
||||
startw = w_max // 2 - w_p // 2
|
||||
endw = startw + w_p
|
||||
original_pe_indexes = original_pe_indexes[starth:endh, startw:endw]
|
||||
return original_pe_indexes.flatten()
|
||||
|
||||
def forward(self, latent):
|
||||
batch_size, num_channels, height, width = latent.size()
|
||||
latent = latent.view(
|
||||
@@ -80,7 +95,8 @@ class AuraFlowPatchEmbed(nn.Module):
|
||||
)
|
||||
latent = latent.permute(0, 2, 4, 1, 3, 5).flatten(-3).flatten(1, 2)
|
||||
latent = self.proj(latent)
|
||||
return latent + self.pos_embed
|
||||
pe_index = self.pe_selection_index_based_on_dim(height, width)
|
||||
return latent + self.pos_embed[:, pe_index]
|
||||
|
||||
|
||||
# Taken from the original Aura flow inference code.
|
||||
|
||||
@@ -13,7 +13,7 @@
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
from typing import Any, Dict, Optional, Union
|
||||
from typing import Any, Dict, Optional, Tuple, Union
|
||||
|
||||
import torch
|
||||
from torch import nn
|
||||
@@ -22,6 +22,7 @@ from ...configuration_utils import ConfigMixin, register_to_config
|
||||
from ...utils import is_torch_version, logging
|
||||
from ...utils.torch_utils import maybe_allow_in_graph
|
||||
from ..attention import Attention, FeedForward
|
||||
from ..attention_processor import AttentionProcessor, CogVideoXAttnProcessor2_0, FusedCogVideoXAttnProcessor2_0
|
||||
from ..embeddings import CogVideoXPatchEmbed, TimestepEmbedding, Timesteps, get_3d_sincos_pos_embed
|
||||
from ..modeling_outputs import Transformer2DModelOutput
|
||||
from ..modeling_utils import ModelMixin
|
||||
@@ -37,13 +38,20 @@ class CogVideoXBlock(nn.Module):
|
||||
Transformer block used in [CogVideoX](https://github.com/THUDM/CogVideo) model.
|
||||
|
||||
Parameters:
|
||||
dim (`int`): The number of channels in the input and output.
|
||||
num_attention_heads (`int`): The number of heads to use for multi-head attention.
|
||||
attention_head_dim (`int`): The number of channels in each head.
|
||||
dropout (`float`, *optional*, defaults to 0.0): The dropout probability to use.
|
||||
activation_fn (`str`, *optional*, defaults to `"geglu"`): Activation function to be used in feed-forward.
|
||||
attention_bias (:
|
||||
obj: `bool`, *optional*, defaults to `False`): Configure if the attentions should contain a bias parameter.
|
||||
dim (`int`):
|
||||
The number of channels in the input and output.
|
||||
num_attention_heads (`int`):
|
||||
The number of heads to use for multi-head attention.
|
||||
attention_head_dim (`int`):
|
||||
The number of channels in each head.
|
||||
time_embed_dim (`int`):
|
||||
The number of channels in timestep embedding.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
The dropout probability to use.
|
||||
activation_fn (`str`, defaults to `"gelu-approximate"`):
|
||||
Activation function to be used in feed-forward.
|
||||
attention_bias (`bool`, defaults to `False`):
|
||||
Whether or not to use bias in attention projection layers.
|
||||
qk_norm (`bool`, defaults to `True`):
|
||||
Whether or not to use normalization after query and key projections in Attention.
|
||||
norm_elementwise_affine (`bool`, defaults to `True`):
|
||||
@@ -90,6 +98,7 @@ class CogVideoXBlock(nn.Module):
|
||||
eps=1e-6,
|
||||
bias=attention_bias,
|
||||
out_bias=attention_out_bias,
|
||||
processor=CogVideoXAttnProcessor2_0(),
|
||||
)
|
||||
|
||||
# 2. Feed Forward
|
||||
@@ -109,24 +118,24 @@ class CogVideoXBlock(nn.Module):
|
||||
hidden_states: torch.Tensor,
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
temb: torch.Tensor,
|
||||
image_rotary_emb: Optional[Tuple[torch.Tensor, torch.Tensor]] = None,
|
||||
) -> torch.Tensor:
|
||||
text_seq_length = encoder_hidden_states.size(1)
|
||||
|
||||
# norm & modulate
|
||||
norm_hidden_states, norm_encoder_hidden_states, gate_msa, enc_gate_msa = self.norm1(
|
||||
hidden_states, encoder_hidden_states, temb
|
||||
)
|
||||
|
||||
# attention
|
||||
text_length = norm_encoder_hidden_states.size(1)
|
||||
|
||||
# CogVideoX uses concatenated text + video embeddings with self-attention instead of using
|
||||
# them in cross-attention individually
|
||||
norm_hidden_states = torch.cat([norm_encoder_hidden_states, norm_hidden_states], dim=1)
|
||||
attn_output = self.attn1(
|
||||
attn_hidden_states, attn_encoder_hidden_states = self.attn1(
|
||||
hidden_states=norm_hidden_states,
|
||||
encoder_hidden_states=None,
|
||||
encoder_hidden_states=norm_encoder_hidden_states,
|
||||
image_rotary_emb=image_rotary_emb,
|
||||
)
|
||||
|
||||
hidden_states = hidden_states + gate_msa * attn_output[:, text_length:]
|
||||
encoder_hidden_states = encoder_hidden_states + enc_gate_msa * attn_output[:, :text_length]
|
||||
hidden_states = hidden_states + gate_msa * attn_hidden_states
|
||||
encoder_hidden_states = encoder_hidden_states + enc_gate_msa * attn_encoder_hidden_states
|
||||
|
||||
# norm & modulate
|
||||
norm_hidden_states, norm_encoder_hidden_states, gate_ff, enc_gate_ff = self.norm2(
|
||||
@@ -137,8 +146,9 @@ class CogVideoXBlock(nn.Module):
|
||||
norm_hidden_states = torch.cat([norm_encoder_hidden_states, norm_hidden_states], dim=1)
|
||||
ff_output = self.ff(norm_hidden_states)
|
||||
|
||||
hidden_states = hidden_states + gate_ff * ff_output[:, text_length:]
|
||||
encoder_hidden_states = encoder_hidden_states + enc_gate_ff * ff_output[:, :text_length]
|
||||
hidden_states = hidden_states + gate_ff * ff_output[:, text_seq_length:]
|
||||
encoder_hidden_states = encoder_hidden_states + enc_gate_ff * ff_output[:, :text_seq_length]
|
||||
|
||||
return hidden_states, encoder_hidden_states
|
||||
|
||||
|
||||
@@ -147,36 +157,53 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
A Transformer model for video-like data in [CogVideoX](https://github.com/THUDM/CogVideo).
|
||||
|
||||
Parameters:
|
||||
num_attention_heads (`int`, *optional*, defaults to 16): The number of heads to use for multi-head attention.
|
||||
attention_head_dim (`int`, *optional*, defaults to 88): The number of channels in each head.
|
||||
in_channels (`int`, *optional*):
|
||||
num_attention_heads (`int`, defaults to `30`):
|
||||
The number of heads to use for multi-head attention.
|
||||
attention_head_dim (`int`, defaults to `64`):
|
||||
The number of channels in each head.
|
||||
in_channels (`int`, defaults to `16`):
|
||||
The number of channels in the input.
|
||||
out_channels (`int`, *optional*):
|
||||
out_channels (`int`, *optional*, defaults to `16`):
|
||||
The number of channels in the output.
|
||||
num_layers (`int`, *optional*, defaults to 1): The number of layers of Transformer blocks to use.
|
||||
dropout (`float`, *optional*, defaults to 0.0): The dropout probability to use.
|
||||
cross_attention_dim (`int`, *optional*): The number of `encoder_hidden_states` dimensions to use.
|
||||
attention_bias (`bool`, *optional*):
|
||||
Configure if the `TransformerBlocks` attention should contain a bias parameter.
|
||||
sample_size (`int`, *optional*): The width of the latent images (specify if the input is **discrete**).
|
||||
This is fixed during training since it is used to learn a number of position embeddings.
|
||||
patch_size (`int`, *optional*):
|
||||
flip_sin_to_cos (`bool`, defaults to `True`):
|
||||
Whether to flip the sin to cos in the time embedding.
|
||||
time_embed_dim (`int`, defaults to `512`):
|
||||
Output dimension of timestep embeddings.
|
||||
text_embed_dim (`int`, defaults to `4096`):
|
||||
Input dimension of text embeddings from the text encoder.
|
||||
num_layers (`int`, defaults to `30`):
|
||||
The number of layers of Transformer blocks to use.
|
||||
dropout (`float`, defaults to `0.0`):
|
||||
The dropout probability to use.
|
||||
attention_bias (`bool`, defaults to `True`):
|
||||
Whether or not to use bias in the attention projection layers.
|
||||
sample_width (`int`, defaults to `90`):
|
||||
The width of the input latents.
|
||||
sample_height (`int`, defaults to `60`):
|
||||
The height of the input latents.
|
||||
sample_frames (`int`, defaults to `49`):
|
||||
The number of frames in the input latents. Note that this parameter was incorrectly initialized to 49
|
||||
instead of 13 because CogVideoX processed 13 latent frames at once in its default and recommended settings,
|
||||
but cannot be changed to the correct value to ensure backwards compatibility. To create a transformer with
|
||||
K latent frames, the correct value to pass here would be: ((K - 1) * temporal_compression_ratio + 1).
|
||||
patch_size (`int`, defaults to `2`):
|
||||
The size of the patches to use in the patch embedding layer.
|
||||
activation_fn (`str`, *optional*, defaults to `"geglu"`): Activation function to use in feed-forward.
|
||||
num_embeds_ada_norm ( `int`, *optional*):
|
||||
The number of diffusion steps used during training. Pass if at least one of the norm_layers is
|
||||
`AdaLayerNorm`. This is fixed during training since it is used to learn a number of embeddings that are
|
||||
added to the hidden states. During inference, you can denoise for up to but not more steps than
|
||||
`num_embeds_ada_norm`.
|
||||
norm_type (`str`, *optional*, defaults to `"layer_norm"`):
|
||||
The type of normalization to use. Options are `"layer_norm"` or `"ada_layer_norm"`.
|
||||
norm_elementwise_affine (`bool`, *optional*, defaults to `True`):
|
||||
temporal_compression_ratio (`int`, defaults to `4`):
|
||||
The compression ratio across the temporal dimension. See documentation for `sample_frames`.
|
||||
max_text_seq_length (`int`, defaults to `226`):
|
||||
The maximum sequence length of the input text embeddings.
|
||||
activation_fn (`str`, defaults to `"gelu-approximate"`):
|
||||
Activation function to use in feed-forward.
|
||||
timestep_activation_fn (`str`, defaults to `"silu"`):
|
||||
Activation function to use when generating the timestep embeddings.
|
||||
norm_elementwise_affine (`bool`, defaults to `True`):
|
||||
Whether or not to use elementwise affine in normalization layers.
|
||||
norm_eps (`float`, *optional*, defaults to 1e-5): The epsilon value to use in normalization layers.
|
||||
caption_channels (`int`, *optional*):
|
||||
The number of channels in the caption embeddings.
|
||||
video_length (`int`, *optional*):
|
||||
The number of frames in the video-like data.
|
||||
norm_eps (`float`, defaults to `1e-5`):
|
||||
The epsilon value to use in normalization layers.
|
||||
spatial_interpolation_scale (`float`, defaults to `1.875`):
|
||||
Scaling factor to apply in 3D positional embeddings across spatial dimensions.
|
||||
temporal_interpolation_scale (`float`, defaults to `1.0`):
|
||||
Scaling factor to apply in 3D positional embeddings across temporal dimensions.
|
||||
"""
|
||||
|
||||
_supports_gradient_checkpointing = True
|
||||
@@ -186,7 +213,7 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
self,
|
||||
num_attention_heads: int = 30,
|
||||
attention_head_dim: int = 64,
|
||||
in_channels: Optional[int] = 16,
|
||||
in_channels: int = 16,
|
||||
out_channels: Optional[int] = 16,
|
||||
flip_sin_to_cos: bool = True,
|
||||
freq_shift: int = 0,
|
||||
@@ -207,6 +234,7 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
norm_eps: float = 1e-5,
|
||||
spatial_interpolation_scale: float = 1.875,
|
||||
temporal_interpolation_scale: float = 1.0,
|
||||
use_rotary_positional_embeddings: bool = False,
|
||||
):
|
||||
super().__init__()
|
||||
inner_dim = num_attention_heads * attention_head_dim
|
||||
@@ -271,12 +299,113 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
def _set_gradient_checkpointing(self, module, value=False):
|
||||
self.gradient_checkpointing = value
|
||||
|
||||
@property
|
||||
# Copied from diffusers.models.unets.unet_2d_condition.UNet2DConditionModel.attn_processors
|
||||
def attn_processors(self) -> Dict[str, AttentionProcessor]:
|
||||
r"""
|
||||
Returns:
|
||||
`dict` of attention processors: A dictionary containing all attention processors used in the model with
|
||||
indexed by its weight name.
|
||||
"""
|
||||
# set recursively
|
||||
processors = {}
|
||||
|
||||
def fn_recursive_add_processors(name: str, module: torch.nn.Module, processors: Dict[str, AttentionProcessor]):
|
||||
if hasattr(module, "get_processor"):
|
||||
processors[f"{name}.processor"] = module.get_processor()
|
||||
|
||||
for sub_name, child in module.named_children():
|
||||
fn_recursive_add_processors(f"{name}.{sub_name}", child, processors)
|
||||
|
||||
return processors
|
||||
|
||||
for name, module in self.named_children():
|
||||
fn_recursive_add_processors(name, module, processors)
|
||||
|
||||
return processors
|
||||
|
||||
# Copied from diffusers.models.unets.unet_2d_condition.UNet2DConditionModel.set_attn_processor
|
||||
def set_attn_processor(self, processor: Union[AttentionProcessor, Dict[str, AttentionProcessor]]):
|
||||
r"""
|
||||
Sets the attention processor to use to compute attention.
|
||||
|
||||
Parameters:
|
||||
processor (`dict` of `AttentionProcessor` or only `AttentionProcessor`):
|
||||
The instantiated processor class or a dictionary of processor classes that will be set as the processor
|
||||
for **all** `Attention` layers.
|
||||
|
||||
If `processor` is a dict, the key needs to define the path to the corresponding cross attention
|
||||
processor. This is strongly recommended when setting trainable attention processors.
|
||||
|
||||
"""
|
||||
count = len(self.attn_processors.keys())
|
||||
|
||||
if isinstance(processor, dict) and len(processor) != count:
|
||||
raise ValueError(
|
||||
f"A dict of processors was passed, but the number of processors {len(processor)} does not match the"
|
||||
f" number of attention layers: {count}. Please make sure to pass {count} processor classes."
|
||||
)
|
||||
|
||||
def fn_recursive_attn_processor(name: str, module: torch.nn.Module, processor):
|
||||
if hasattr(module, "set_processor"):
|
||||
if not isinstance(processor, dict):
|
||||
module.set_processor(processor)
|
||||
else:
|
||||
module.set_processor(processor.pop(f"{name}.processor"))
|
||||
|
||||
for sub_name, child in module.named_children():
|
||||
fn_recursive_attn_processor(f"{name}.{sub_name}", child, processor)
|
||||
|
||||
for name, module in self.named_children():
|
||||
fn_recursive_attn_processor(name, module, processor)
|
||||
|
||||
# Copied from diffusers.models.unets.unet_2d_condition.UNet2DConditionModel.fuse_qkv_projections with FusedAttnProcessor2_0->FusedCogVideoXAttnProcessor2_0
|
||||
def fuse_qkv_projections(self):
|
||||
"""
|
||||
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query, key, value)
|
||||
are fused. For cross-attention modules, key and value projection matrices are fused.
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
This API is 🧪 experimental.
|
||||
|
||||
</Tip>
|
||||
"""
|
||||
self.original_attn_processors = None
|
||||
|
||||
for _, attn_processor in self.attn_processors.items():
|
||||
if "Added" in str(attn_processor.__class__.__name__):
|
||||
raise ValueError("`fuse_qkv_projections()` is not supported for models having added KV projections.")
|
||||
|
||||
self.original_attn_processors = self.attn_processors
|
||||
|
||||
for module in self.modules():
|
||||
if isinstance(module, Attention):
|
||||
module.fuse_projections(fuse=True)
|
||||
|
||||
self.set_attn_processor(FusedCogVideoXAttnProcessor2_0())
|
||||
|
||||
# Copied from diffusers.models.unets.unet_2d_condition.UNet2DConditionModel.unfuse_qkv_projections
|
||||
def unfuse_qkv_projections(self):
|
||||
"""Disables the fused QKV projection if enabled.
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
This API is 🧪 experimental.
|
||||
|
||||
</Tip>
|
||||
|
||||
"""
|
||||
if self.original_attn_processors is not None:
|
||||
self.set_attn_processor(self.original_attn_processors)
|
||||
|
||||
def forward(
|
||||
self,
|
||||
hidden_states: torch.Tensor,
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
timestep: Union[int, float, torch.LongTensor],
|
||||
timestep_cond: Optional[torch.Tensor] = None,
|
||||
image_rotary_emb: Optional[Tuple[torch.Tensor, torch.Tensor]] = None,
|
||||
return_dict: bool = True,
|
||||
):
|
||||
batch_size, num_frames, channels, height, width = hidden_states.shape
|
||||
@@ -295,16 +424,18 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
hidden_states = self.patch_embed(encoder_hidden_states, hidden_states)
|
||||
|
||||
# 3. Position embedding
|
||||
seq_length = height * width * num_frames // (self.config.patch_size**2)
|
||||
text_seq_length = encoder_hidden_states.shape[1]
|
||||
if not self.config.use_rotary_positional_embeddings:
|
||||
seq_length = height * width * num_frames // (self.config.patch_size**2)
|
||||
|
||||
pos_embeds = self.pos_embedding[:, : self.config.max_text_seq_length + seq_length]
|
||||
hidden_states = hidden_states + pos_embeds
|
||||
hidden_states = self.embedding_dropout(hidden_states)
|
||||
pos_embeds = self.pos_embedding[:, : text_seq_length + seq_length]
|
||||
hidden_states = hidden_states + pos_embeds
|
||||
hidden_states = self.embedding_dropout(hidden_states)
|
||||
|
||||
encoder_hidden_states = hidden_states[:, : self.config.max_text_seq_length]
|
||||
hidden_states = hidden_states[:, self.config.max_text_seq_length :]
|
||||
encoder_hidden_states = hidden_states[:, :text_seq_length]
|
||||
hidden_states = hidden_states[:, text_seq_length:]
|
||||
|
||||
# 5. Transformer blocks
|
||||
# 4. Transformer blocks
|
||||
for i, block in enumerate(self.transformer_blocks):
|
||||
if self.training and self.gradient_checkpointing:
|
||||
|
||||
@@ -320,6 +451,7 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
hidden_states,
|
||||
encoder_hidden_states,
|
||||
emb,
|
||||
image_rotary_emb,
|
||||
**ckpt_kwargs,
|
||||
)
|
||||
else:
|
||||
@@ -327,15 +459,23 @@ class CogVideoXTransformer3DModel(ModelMixin, ConfigMixin):
|
||||
hidden_states=hidden_states,
|
||||
encoder_hidden_states=encoder_hidden_states,
|
||||
temb=emb,
|
||||
image_rotary_emb=image_rotary_emb,
|
||||
)
|
||||
|
||||
hidden_states = self.norm_final(hidden_states)
|
||||
if not self.config.use_rotary_positional_embeddings:
|
||||
# CogVideoX-2B
|
||||
hidden_states = self.norm_final(hidden_states)
|
||||
else:
|
||||
# CogVideoX-5B
|
||||
hidden_states = torch.cat([encoder_hidden_states, hidden_states], dim=1)
|
||||
hidden_states = self.norm_final(hidden_states)
|
||||
hidden_states = hidden_states[:, text_seq_length:]
|
||||
|
||||
# 6. Final block
|
||||
# 5. Final block
|
||||
hidden_states = self.norm_out(hidden_states, temb=emb)
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
|
||||
# 7. Unpatchify
|
||||
# 6. Unpatchify
|
||||
p = self.config.patch_size
|
||||
output = hidden_states.reshape(batch_size, num_frames, height // p, width // p, channels, p, p)
|
||||
output = output.permute(0, 1, 4, 2, 5, 3, 6).flatten(5, 6).flatten(3, 4)
|
||||
|
||||
@@ -23,6 +23,7 @@ from transformers import T5EncoderModel, T5Tokenizer
|
||||
|
||||
from ...callbacks import MultiPipelineCallbacks, PipelineCallback
|
||||
from ...models import AutoencoderKLCogVideoX, CogVideoXTransformer3DModel
|
||||
from ...models.embeddings import get_3d_rotary_pos_embed
|
||||
from ...pipelines.pipeline_utils import DiffusionPipeline
|
||||
from ...schedulers import CogVideoXDDIMScheduler, CogVideoXDPMScheduler
|
||||
from ...utils import BaseOutput, logging, replace_example_docstring
|
||||
@@ -40,6 +41,7 @@ EXAMPLE_DOC_STRING = """
|
||||
>>> from diffusers import CogVideoXPipeline
|
||||
>>> from diffusers.utils import export_to_video
|
||||
|
||||
>>> # Models: "THUDM/CogVideoX-2b" or "THUDM/CogVideoX-5b"
|
||||
>>> pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-2b", torch_dtype=torch.float16).to("cuda")
|
||||
>>> prompt = (
|
||||
... "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
|
||||
@@ -55,6 +57,25 @@ EXAMPLE_DOC_STRING = """
|
||||
"""
|
||||
|
||||
|
||||
# Similar to diffusers.pipelines.hunyuandit.pipeline_hunyuandit.get_resize_crop_region_for_grid
|
||||
def get_resize_crop_region_for_grid(src, tgt_width, tgt_height):
|
||||
tw = tgt_width
|
||||
th = tgt_height
|
||||
h, w = src
|
||||
r = h / w
|
||||
if r > (th / tw):
|
||||
resize_height = th
|
||||
resize_width = int(round(th / h * w))
|
||||
else:
|
||||
resize_width = tw
|
||||
resize_height = int(round(tw / w * h))
|
||||
|
||||
crop_top = int(round((th - resize_height) / 2.0))
|
||||
crop_left = int(round((tw - resize_width) / 2.0))
|
||||
|
||||
return (crop_top, crop_left), (crop_top + resize_height, crop_left + resize_width)
|
||||
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.retrieve_timesteps
|
||||
def retrieve_timesteps(
|
||||
scheduler,
|
||||
@@ -332,20 +353,11 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
latents = latents * self.scheduler.init_noise_sigma
|
||||
return latents
|
||||
|
||||
def decode_latents(self, latents: torch.Tensor, num_seconds: int):
|
||||
def decode_latents(self, latents: torch.Tensor) -> torch.Tensor:
|
||||
latents = latents.permute(0, 2, 1, 3, 4) # [batch_size, num_channels, num_frames, height, width]
|
||||
latents = 1 / self.vae.config.scaling_factor * latents
|
||||
|
||||
frames = []
|
||||
for i in range(num_seconds):
|
||||
start_frame, end_frame = (0, 3) if i == 0 else (2 * i + 1, 2 * i + 3)
|
||||
|
||||
current_frames = self.vae.decode(latents[:, :, start_frame:end_frame]).sample
|
||||
frames.append(current_frames)
|
||||
|
||||
self.vae.clear_fake_context_parallel_cache()
|
||||
|
||||
frames = torch.cat(frames, dim=2)
|
||||
frames = self.vae.decode(latents).sample
|
||||
return frames
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs
|
||||
@@ -418,6 +430,46 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
f" {negative_prompt_embeds.shape}."
|
||||
)
|
||||
|
||||
def fuse_qkv_projections(self) -> None:
|
||||
r"""Enables fused QKV projections."""
|
||||
self.fusing_transformer = True
|
||||
self.transformer.fuse_qkv_projections()
|
||||
|
||||
def unfuse_qkv_projections(self) -> None:
|
||||
r"""Disable QKV projection fusion if enabled."""
|
||||
if not self.fusing_transformer:
|
||||
logger.warning("The Transformer was not initially fused for QKV projections. Doing nothing.")
|
||||
else:
|
||||
self.transformer.unfuse_qkv_projections()
|
||||
self.fusing_transformer = False
|
||||
|
||||
def _prepare_rotary_positional_embeddings(
|
||||
self,
|
||||
height: int,
|
||||
width: int,
|
||||
num_frames: int,
|
||||
device: torch.device,
|
||||
) -> Tuple[torch.Tensor, torch.Tensor]:
|
||||
grid_height = height // (self.vae_scale_factor_spatial * self.transformer.config.patch_size)
|
||||
grid_width = width // (self.vae_scale_factor_spatial * self.transformer.config.patch_size)
|
||||
base_size_width = 720 // (self.vae_scale_factor_spatial * self.transformer.config.patch_size)
|
||||
base_size_height = 480 // (self.vae_scale_factor_spatial * self.transformer.config.patch_size)
|
||||
|
||||
grid_crops_coords = get_resize_crop_region_for_grid(
|
||||
(grid_height, grid_width), base_size_width, base_size_height
|
||||
)
|
||||
freqs_cos, freqs_sin = get_3d_rotary_pos_embed(
|
||||
embed_dim=self.transformer.config.attention_head_dim,
|
||||
crops_coords=grid_crops_coords,
|
||||
grid_size=(grid_height, grid_width),
|
||||
temporal_size=num_frames,
|
||||
use_real=True,
|
||||
)
|
||||
|
||||
freqs_cos = freqs_cos.to(device=device)
|
||||
freqs_sin = freqs_sin.to(device=device)
|
||||
return freqs_cos, freqs_sin
|
||||
|
||||
@property
|
||||
def guidance_scale(self):
|
||||
return self._guidance_scale
|
||||
@@ -438,8 +490,7 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
height: int = 480,
|
||||
width: int = 720,
|
||||
num_frames: int = 48,
|
||||
fps: int = 8,
|
||||
num_frames: int = 49,
|
||||
num_inference_steps: int = 50,
|
||||
timesteps: Optional[List[int]] = None,
|
||||
guidance_scale: float = 6,
|
||||
@@ -534,9 +585,10 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
`tuple`. When returning a tuple, the first element is a list with the generated images.
|
||||
"""
|
||||
|
||||
assert (
|
||||
num_frames <= 48 and num_frames % fps == 0 and fps == 8
|
||||
), f"The number of frames must be divisible by {fps=} and less than 48 frames (for now). Other values are not supported in CogVideoX."
|
||||
if num_frames > 49:
|
||||
raise ValueError(
|
||||
"The number of frames must be less than 49 for now due to static positional embeddings. This will be updated in the future to remove this limitation."
|
||||
)
|
||||
|
||||
if isinstance(callback_on_step_end, (PipelineCallback, MultiPipelineCallbacks)):
|
||||
callback_on_step_end_tensor_inputs = callback_on_step_end.tensor_inputs
|
||||
@@ -593,7 +645,6 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
|
||||
# 5. Prepare latents.
|
||||
latent_channels = self.transformer.config.in_channels
|
||||
num_frames += 1
|
||||
latents = self.prepare_latents(
|
||||
batch_size * num_videos_per_prompt,
|
||||
latent_channels,
|
||||
@@ -609,7 +660,14 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 7. Denoising loop
|
||||
# 7. Create rotary embeds if required
|
||||
image_rotary_emb = (
|
||||
self._prepare_rotary_positional_embeddings(height, width, latents.size(1), device)
|
||||
if self.transformer.config.use_rotary_positional_embeddings
|
||||
else None
|
||||
)
|
||||
|
||||
# 8. Denoising loop
|
||||
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0)
|
||||
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
@@ -630,6 +688,7 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
hidden_states=latent_model_input,
|
||||
encoder_hidden_states=prompt_embeds,
|
||||
timestep=timestep,
|
||||
image_rotary_emb=image_rotary_emb,
|
||||
return_dict=False,
|
||||
)[0]
|
||||
noise_pred = noise_pred.float()
|
||||
@@ -673,7 +732,7 @@ class CogVideoXPipeline(DiffusionPipeline):
|
||||
progress_bar.update()
|
||||
|
||||
if not output_type == "latent":
|
||||
video = self.decode_latents(latents, num_frames // fps)
|
||||
video = self.decode_latents(latents)
|
||||
video = self.video_processor.postprocess_video(video=video, output_type=output_type)
|
||||
else:
|
||||
video = latents
|
||||
|
||||
@@ -280,7 +280,7 @@ class FluxPipeline(DiffusionPipeline, FluxLoraLoaderMixin):
|
||||
prompt_embeds = prompt_embeds.to(dtype=self.text_encoder.dtype, device=device)
|
||||
|
||||
# duplicate text embeddings for each generation per prompt, using mps friendly method
|
||||
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
|
||||
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt)
|
||||
prompt_embeds = prompt_embeds.view(batch_size * num_images_per_prompt, -1)
|
||||
|
||||
return prompt_embeds
|
||||
|
||||
@@ -19,7 +19,6 @@ import inspect
|
||||
import os
|
||||
import re
|
||||
import sys
|
||||
from collections import OrderedDict
|
||||
from dataclasses import dataclass
|
||||
from pathlib import Path
|
||||
from typing import Any, Callable, Dict, List, Optional, Union, get_args, get_origin
|
||||
@@ -1173,93 +1172,6 @@ class DiffusionPipeline(ConfigMixin, PushToHubMixin):
|
||||
component.to("cpu")
|
||||
self.hf_device_map = None
|
||||
|
||||
def enable_dynamic_upcasting(
|
||||
self,
|
||||
components: Optional[List[str]] = None,
|
||||
upcast_dtype: Optional[torch.dtype] = None,
|
||||
):
|
||||
r"""
|
||||
Enable module-wise dynamic upcasting. This allows models to be loaded into the GPU in a low memory dtype e.g.
|
||||
torch.float8_e4m3fn, but perform inference using a dtype that is supported on the GPU, by casting the module to
|
||||
the appropriate dtype right before the foward pass. The module is then moved back to the low memory dtype after
|
||||
the foward pass.
|
||||
|
||||
"""
|
||||
if components is None:
|
||||
raise ValueError("Please provide a list of pipeline component names to apply dynamic upcasting")
|
||||
|
||||
def fn_recursive_upcast(module, dtype, original_dtype, keep_in_fp32_modules):
|
||||
has_children = list(module.children())
|
||||
upcast_dtype = dtype
|
||||
downcast_dtype = original_dtype
|
||||
|
||||
def upcast_hook_fn(module, inputs):
|
||||
module = module.to(upcast_dtype)
|
||||
|
||||
def downcast_hook_fn(module, *args, **kwargs):
|
||||
module = module.to(downcast_dtype)
|
||||
|
||||
if not has_children:
|
||||
module.register_forward_pre_hook(upcast_hook_fn)
|
||||
module.register_forward_hook(downcast_hook_fn)
|
||||
|
||||
for name, child in module.named_children():
|
||||
if any(module_to_keep_in_fp32 in name.split(".") for module_to_keep_in_fp32 in keep_in_fp32_modules):
|
||||
dtype = torch.float32
|
||||
else:
|
||||
dtype = upcast_dtype
|
||||
|
||||
fn_recursive_upcast(child, dtype, original_dtype, keep_in_fp32_modules)
|
||||
|
||||
for component in components:
|
||||
if not hasattr(self, component):
|
||||
raise ValueError(f"Pipeline has no component named: {component}")
|
||||
|
||||
component_module = getattr(self, component)
|
||||
if not isinstance(component_module, torch.nn.Module):
|
||||
raise ValueError(
|
||||
f"Pipeline component: {component} is not a torch.nn.Module. Cannot apply dynamic upcasting."
|
||||
)
|
||||
|
||||
use_keep_in_fp32_modules = (
|
||||
hasattr(component_module, "_keep_in_fp32_modules")
|
||||
and (component_module._keep_in_fp32_modules is not None)
|
||||
and (upcast_dtype != torch.float32)
|
||||
)
|
||||
if use_keep_in_fp32_modules:
|
||||
keep_in_fp32_modules = component_module._keep_in_fp32_modules
|
||||
else:
|
||||
keep_in_fp32_modules = []
|
||||
|
||||
original_dtype = component_module.dtype
|
||||
for name, module in component_module.named_children():
|
||||
fn_recursive_upcast(module, upcast_dtype, original_dtype, keep_in_fp32_modules)
|
||||
|
||||
def disable_dynamic_upcasting(
|
||||
self,
|
||||
):
|
||||
def fn_recursive_upcast(module):
|
||||
has_children = list(module.children())
|
||||
if not has_children:
|
||||
module._forward_pre_hooks = OrderedDict()
|
||||
module._forward_hooks = OrderedDict()
|
||||
|
||||
for child in module.children():
|
||||
fn_recursive_upcast(child)
|
||||
|
||||
for component in self.components:
|
||||
if not hasattr(self, component):
|
||||
raise ValueError(f"Pipeline has no component named: {component}")
|
||||
|
||||
component_module = getattr(self, component)
|
||||
if not issubclass(component_module, torch.nn.Module):
|
||||
raise ValueError(
|
||||
f"Pipeline component: {component} is not an torch.nn.Module. Cannot apply dynamic upcasting."
|
||||
)
|
||||
|
||||
for module in component_module.children():
|
||||
fn_recursive_upcast(module)
|
||||
|
||||
@classmethod
|
||||
@validate_hf_hub_args
|
||||
def download(cls, pretrained_model_name, **kwargs) -> Union[str, os.PathLike]:
|
||||
|
||||
@@ -12,19 +12,26 @@
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
import os
|
||||
import sys
|
||||
import tempfile
|
||||
import unittest
|
||||
|
||||
import numpy as np
|
||||
import safetensors.torch
|
||||
import torch
|
||||
from transformers import AutoTokenizer, CLIPTextModel, CLIPTokenizer, T5EncoderModel
|
||||
|
||||
from diffusers import FlowMatchEulerDiscreteScheduler, FluxPipeline, FluxTransformer2DModel
|
||||
from diffusers.utils.testing_utils import floats_tensor, require_peft_backend
|
||||
from diffusers.utils.testing_utils import floats_tensor, is_peft_available, require_peft_backend, torch_device
|
||||
|
||||
|
||||
if is_peft_available():
|
||||
from peft.utils import get_peft_model_state_dict
|
||||
|
||||
sys.path.append(".")
|
||||
|
||||
from utils import PeftLoraLoaderMixinTests # noqa: E402
|
||||
from utils import PeftLoraLoaderMixinTests, check_if_lora_correctly_set # noqa: E402
|
||||
|
||||
|
||||
@require_peft_backend
|
||||
@@ -90,3 +97,51 @@ class FluxLoRATests(unittest.TestCase, PeftLoraLoaderMixinTests):
|
||||
pipeline_inputs.update({"generator": generator})
|
||||
|
||||
return noise, input_ids, pipeline_inputs
|
||||
|
||||
def test_with_alpha_in_state_dict(self):
|
||||
components, _, denoiser_lora_config = self.get_dummy_components(FlowMatchEulerDiscreteScheduler)
|
||||
pipe = self.pipeline_class(**components)
|
||||
pipe = pipe.to(torch_device)
|
||||
pipe.set_progress_bar_config(disable=None)
|
||||
_, _, inputs = self.get_dummy_inputs(with_generator=False)
|
||||
|
||||
output_no_lora = pipe(**inputs, generator=torch.manual_seed(0)).images
|
||||
self.assertTrue(output_no_lora.shape == self.output_shape)
|
||||
|
||||
pipe.transformer.add_adapter(denoiser_lora_config)
|
||||
self.assertTrue(check_if_lora_correctly_set(pipe.transformer), "Lora not correctly set in transformer")
|
||||
|
||||
images_lora = pipe(**inputs, generator=torch.manual_seed(0)).images
|
||||
|
||||
with tempfile.TemporaryDirectory() as tmpdirname:
|
||||
denoiser_state_dict = get_peft_model_state_dict(pipe.transformer)
|
||||
self.pipeline_class.save_lora_weights(tmpdirname, transformer_lora_layers=denoiser_state_dict)
|
||||
|
||||
self.assertTrue(os.path.isfile(os.path.join(tmpdirname, "pytorch_lora_weights.safetensors")))
|
||||
pipe.unload_lora_weights()
|
||||
pipe.load_lora_weights(os.path.join(tmpdirname, "pytorch_lora_weights.safetensors"))
|
||||
|
||||
# modify the state dict to have alpha values following
|
||||
# https://huggingface.co/TheLastBen/Jon_Snow_Flux_LoRA/blob/main/jon_snow.safetensors
|
||||
state_dict_with_alpha = safetensors.torch.load_file(
|
||||
os.path.join(tmpdirname, "pytorch_lora_weights.safetensors")
|
||||
)
|
||||
alpha_dict = {}
|
||||
for k, v in state_dict_with_alpha.items():
|
||||
# only do for `transformer` and for the k projections -- should be enough to test.
|
||||
if "transformer" in k and "to_k" in k and "lora_A" in k:
|
||||
alpha_dict[f"{k}.alpha"] = float(torch.randint(10, 100, size=()))
|
||||
state_dict_with_alpha.update(alpha_dict)
|
||||
|
||||
images_lora_from_pretrained = pipe(**inputs, generator=torch.manual_seed(0)).images
|
||||
self.assertTrue(check_if_lora_correctly_set(pipe.transformer), "Lora not correctly set in denoiser")
|
||||
|
||||
pipe.unload_lora_weights()
|
||||
pipe.load_lora_weights(state_dict_with_alpha)
|
||||
images_lora_with_alpha = pipe(**inputs, generator=torch.manual_seed(0)).images
|
||||
|
||||
self.assertTrue(
|
||||
np.allclose(images_lora, images_lora_from_pretrained, atol=1e-3, rtol=1e-3),
|
||||
"Loading from saved checkpoints should give same results.",
|
||||
)
|
||||
self.assertFalse(np.allclose(images_lora_with_alpha, images_lora, atol=1e-3, rtol=1e-3))
|
||||
|
||||
@@ -0,0 +1,82 @@
|
||||
# coding=utf-8
|
||||
# Copyright 2024 HuggingFace Inc.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
import unittest
|
||||
|
||||
import torch
|
||||
|
||||
from diffusers import CogVideoXTransformer3DModel
|
||||
from diffusers.utils.testing_utils import (
|
||||
enable_full_determinism,
|
||||
torch_device,
|
||||
)
|
||||
|
||||
from ..test_modeling_common import ModelTesterMixin
|
||||
|
||||
|
||||
enable_full_determinism()
|
||||
|
||||
|
||||
class CogVideoXTransformerTests(ModelTesterMixin, unittest.TestCase):
|
||||
model_class = CogVideoXTransformer3DModel
|
||||
main_input_name = "hidden_states"
|
||||
|
||||
@property
|
||||
def dummy_input(self):
|
||||
batch_size = 2
|
||||
num_channels = 4
|
||||
num_frames = 1
|
||||
height = 8
|
||||
width = 8
|
||||
embedding_dim = 8
|
||||
sequence_length = 8
|
||||
|
||||
hidden_states = torch.randn((batch_size, num_frames, num_channels, height, width)).to(torch_device)
|
||||
encoder_hidden_states = torch.randn((batch_size, sequence_length, embedding_dim)).to(torch_device)
|
||||
timestep = torch.randint(0, 1000, size=(batch_size,)).to(torch_device)
|
||||
|
||||
return {
|
||||
"hidden_states": hidden_states,
|
||||
"encoder_hidden_states": encoder_hidden_states,
|
||||
"timestep": timestep,
|
||||
}
|
||||
|
||||
@property
|
||||
def input_shape(self):
|
||||
return (1, 4, 8, 8)
|
||||
|
||||
@property
|
||||
def output_shape(self):
|
||||
return (1, 4, 8, 8)
|
||||
|
||||
def prepare_init_args_and_inputs_for_common(self):
|
||||
init_dict = {
|
||||
# Product of num_attention_heads * attention_head_dim must be divisible by 16 for 3D positional embeddings.
|
||||
"num_attention_heads": 2,
|
||||
"attention_head_dim": 8,
|
||||
"in_channels": 4,
|
||||
"out_channels": 4,
|
||||
"time_embed_dim": 2,
|
||||
"text_embed_dim": 8,
|
||||
"num_layers": 1,
|
||||
"sample_width": 8,
|
||||
"sample_height": 8,
|
||||
"sample_frames": 8,
|
||||
"patch_size": 2,
|
||||
"temporal_compression_ratio": 4,
|
||||
"max_text_seq_length": 8,
|
||||
}
|
||||
inputs_dict = self.dummy_input
|
||||
return init_dict, inputs_dict
|
||||
@@ -20,6 +20,7 @@ from diffusers import (
|
||||
)
|
||||
from diffusers.models.attention import FreeNoiseTransformerBlock
|
||||
from diffusers.utils import logging
|
||||
from diffusers.utils.import_utils import is_xformers_available
|
||||
from diffusers.utils.testing_utils import torch_device
|
||||
|
||||
from ..pipeline_params import TEXT_TO_IMAGE_BATCH_PARAMS, TEXT_TO_IMAGE_PARAMS
|
||||
@@ -329,6 +330,13 @@ class AnimateDiffControlNetPipelineFastTests(
|
||||
inputs["prompt_embeds"] = torch.randn((1, 4, pipe.text_encoder.config.hidden_size), device=torch_device)
|
||||
pipe(**inputs)
|
||||
|
||||
@unittest.skipIf(
|
||||
torch_device != "cuda" or not is_xformers_available(),
|
||||
reason="XFormers attention is only available with CUDA and `xformers` installed",
|
||||
)
|
||||
def test_xformers_attention_forwardGenerator_pass(self):
|
||||
super()._test_xformers_attention_forwardGenerator_pass(test_mean_pixel_difference=False)
|
||||
|
||||
def test_free_init(self):
|
||||
components = self.get_dummy_components()
|
||||
pipe: AnimateDiffControlNetPipeline = self.pipeline_class(**components)
|
||||
|
||||
@@ -19,6 +19,7 @@ from diffusers import (
|
||||
UNetMotionModel,
|
||||
)
|
||||
from diffusers.utils import logging
|
||||
from diffusers.utils.import_utils import is_xformers_available
|
||||
from diffusers.utils.testing_utils import torch_device
|
||||
|
||||
from ..pipeline_params import TEXT_TO_IMAGE_BATCH_PARAMS, TEXT_TO_IMAGE_PARAMS
|
||||
@@ -393,6 +394,13 @@ class AnimateDiffSparseControlNetPipelineFastTests(
|
||||
inputs["prompt_embeds"] = torch.randn((1, 4, pipe.text_encoder.config.hidden_size), device=torch_device)
|
||||
pipe(**inputs)
|
||||
|
||||
@unittest.skipIf(
|
||||
torch_device != "cuda" or not is_xformers_available(),
|
||||
reason="XFormers attention is only available with CUDA and `xformers` installed",
|
||||
)
|
||||
def test_xformers_attention_forwardGenerator_pass(self):
|
||||
super()._test_xformers_attention_forwardGenerator_pass(test_mean_pixel_difference=False)
|
||||
|
||||
def test_free_init(self):
|
||||
components = self.get_dummy_components()
|
||||
pipe: AnimateDiffSparseControlNetPipeline = self.pipeline_class(**components)
|
||||
|
||||
@@ -30,7 +30,12 @@ from diffusers.utils.testing_utils import (
|
||||
)
|
||||
|
||||
from ..pipeline_params import TEXT_TO_IMAGE_BATCH_PARAMS, TEXT_TO_IMAGE_IMAGE_PARAMS, TEXT_TO_IMAGE_PARAMS
|
||||
from ..test_pipelines_common import PipelineTesterMixin, to_np
|
||||
from ..test_pipelines_common import (
|
||||
PipelineTesterMixin,
|
||||
check_qkv_fusion_matches_attn_procs_length,
|
||||
check_qkv_fusion_processors_exist,
|
||||
to_np,
|
||||
)
|
||||
|
||||
|
||||
enable_full_determinism()
|
||||
@@ -125,11 +130,6 @@ class CogVideoXPipelineFastTests(PipelineTesterMixin, unittest.TestCase):
|
||||
# Cannot reduce because convolution kernel becomes bigger than sample
|
||||
"height": 16,
|
||||
"width": 16,
|
||||
# TODO(aryan): improve this
|
||||
# Cannot make this lower due to assert condition in pipeline at the moment.
|
||||
# The reason why 8 can't be used here is due to how context-parallel cache works where the first
|
||||
# second of video is decoded from latent frames (0, 3) instead of [(0, 2), (2, 3)]. If 8 is used,
|
||||
# the number of output frames that you get are 5.
|
||||
"num_frames": 8,
|
||||
"max_sequence_length": 16,
|
||||
"output_type": "pt",
|
||||
@@ -148,8 +148,8 @@ class CogVideoXPipelineFastTests(PipelineTesterMixin, unittest.TestCase):
|
||||
video = pipe(**inputs).frames
|
||||
generated_video = video[0]
|
||||
|
||||
self.assertEqual(generated_video.shape, (9, 3, 16, 16))
|
||||
expected_video = torch.randn(9, 3, 16, 16)
|
||||
self.assertEqual(generated_video.shape, (8, 3, 16, 16))
|
||||
expected_video = torch.randn(8, 3, 16, 16)
|
||||
max_diff = np.abs(generated_video - expected_video).max()
|
||||
self.assertLessEqual(max_diff, 1e10)
|
||||
|
||||
@@ -250,6 +250,78 @@ class CogVideoXPipelineFastTests(PipelineTesterMixin, unittest.TestCase):
|
||||
"Attention slicing should not affect the inference results",
|
||||
)
|
||||
|
||||
def test_vae_tiling(self, expected_diff_max: float = 0.2):
|
||||
generator_device = "cpu"
|
||||
components = self.get_dummy_components()
|
||||
|
||||
pipe = self.pipeline_class(**components)
|
||||
pipe.to("cpu")
|
||||
pipe.set_progress_bar_config(disable=None)
|
||||
|
||||
# Without tiling
|
||||
inputs = self.get_dummy_inputs(generator_device)
|
||||
inputs["height"] = inputs["width"] = 128
|
||||
output_without_tiling = pipe(**inputs)[0]
|
||||
|
||||
# With tiling
|
||||
pipe.vae.enable_tiling(
|
||||
tile_sample_min_height=96,
|
||||
tile_sample_min_width=96,
|
||||
tile_overlap_factor_height=1 / 12,
|
||||
tile_overlap_factor_width=1 / 12,
|
||||
)
|
||||
inputs = self.get_dummy_inputs(generator_device)
|
||||
inputs["height"] = inputs["width"] = 128
|
||||
output_with_tiling = pipe(**inputs)[0]
|
||||
|
||||
self.assertLess(
|
||||
(to_np(output_without_tiling) - to_np(output_with_tiling)).max(),
|
||||
expected_diff_max,
|
||||
"VAE tiling should not affect the inference results",
|
||||
)
|
||||
|
||||
@unittest.skip("xformers attention processor does not exist for CogVideoX")
|
||||
def test_xformers_attention_forwardGenerator_pass(self):
|
||||
pass
|
||||
|
||||
def test_fused_qkv_projections(self):
|
||||
device = "cpu" # ensure determinism for the device-dependent torch.Generator
|
||||
components = self.get_dummy_components()
|
||||
pipe = self.pipeline_class(**components)
|
||||
pipe = pipe.to(device)
|
||||
pipe.set_progress_bar_config(disable=None)
|
||||
|
||||
inputs = self.get_dummy_inputs(device)
|
||||
frames = pipe(**inputs).frames # [B, F, C, H, W]
|
||||
original_image_slice = frames[0, -2:, -1, -3:, -3:]
|
||||
|
||||
pipe.fuse_qkv_projections()
|
||||
assert check_qkv_fusion_processors_exist(
|
||||
pipe.transformer
|
||||
), "Something wrong with the fused attention processors. Expected all the attention processors to be fused."
|
||||
assert check_qkv_fusion_matches_attn_procs_length(
|
||||
pipe.transformer, pipe.transformer.original_attn_processors
|
||||
), "Something wrong with the attention processors concerning the fused QKV projections."
|
||||
|
||||
inputs = self.get_dummy_inputs(device)
|
||||
frames = pipe(**inputs).frames
|
||||
image_slice_fused = frames[0, -2:, -1, -3:, -3:]
|
||||
|
||||
pipe.transformer.unfuse_qkv_projections()
|
||||
inputs = self.get_dummy_inputs(device)
|
||||
frames = pipe(**inputs).frames
|
||||
image_slice_disabled = frames[0, -2:, -1, -3:, -3:]
|
||||
|
||||
assert np.allclose(
|
||||
original_image_slice, image_slice_fused, atol=1e-3, rtol=1e-3
|
||||
), "Fusion of QKV projections shouldn't affect the outputs."
|
||||
assert np.allclose(
|
||||
image_slice_fused, image_slice_disabled, atol=1e-3, rtol=1e-3
|
||||
), "Outputs, with QKV projection fusion enabled, shouldn't change when fused QKV projections are disabled."
|
||||
assert np.allclose(
|
||||
original_image_slice, image_slice_disabled, atol=1e-2, rtol=1e-2
|
||||
), "Original outputs should match when fused QKV projections are disabled."
|
||||
|
||||
|
||||
@slow
|
||||
@require_torch_gpu
|
||||
|
||||
@@ -28,6 +28,7 @@ from diffusers import (
|
||||
LattePipeline,
|
||||
LatteTransformer3DModel,
|
||||
)
|
||||
from diffusers.utils.import_utils import is_xformers_available
|
||||
from diffusers.utils.testing_utils import (
|
||||
enable_full_determinism,
|
||||
numpy_cosine_similarity_distance,
|
||||
@@ -256,6 +257,13 @@ class LattePipelineFastTests(PipelineTesterMixin, unittest.TestCase):
|
||||
max_diff = np.abs(to_np(output) - to_np(output_loaded)).max()
|
||||
self.assertLess(max_diff, 1.0)
|
||||
|
||||
@unittest.skipIf(
|
||||
torch_device != "cuda" or not is_xformers_available(),
|
||||
reason="XFormers attention is only available with CUDA and `xformers` installed",
|
||||
)
|
||||
def test_xformers_attention_forwardGenerator_pass(self):
|
||||
super()._test_xformers_attention_forwardGenerator_pass(test_mean_pixel_difference=False)
|
||||
|
||||
|
||||
@slow
|
||||
@require_torch_gpu
|
||||
|
||||
Reference in New Issue
Block a user