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3 Commits

Author SHA1 Message Date
Wauplin be139feda5 fix 2023-12-06 11:34:17 +01:00
Wauplin 6b7716d358 fix import 2023-12-05 19:33:11 +01:00
Wauplin c1b2289529 Harmonize HF environment variables + deprecate use_auth_token 2023-12-05 19:07:04 +01:00
64 changed files with 397 additions and 5955 deletions
+7 -1
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@@ -1,6 +1,12 @@
name: Fast tests for PRs - Test Fetcher
on: workflow_dispatch
on:
pull_request:
branches:
- main
push:
branches:
- ci-*
env:
DIFFUSERS_IS_CI: yes
-1
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@@ -113,7 +113,6 @@ jobs:
- name: Run example PyTorch CPU tests
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples
-4
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@@ -264,10 +264,6 @@
title: ControlNet
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/cycle_diffusion
title: Cycle Diffusion
- local: api/pipelines/dance_diffusion
-3
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@@ -20,9 +20,6 @@ An attention processor is a class for applying different types of attention mech
## AttnProcessor2_0
[[autodoc]] models.attention_processor.AttnProcessor2_0
## FusedAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedAttnProcessor2_0
## LoRAAttnProcessor
[[autodoc]] models.attention_processor.LoRAAttnProcessor
@@ -1,39 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
ControlNet-XS generates images with comparable quality to a regular ControlNet, but it is 20-25% faster ([see benchmark](https://github.com/UmerHA/controlnet-xs-benchmark/blob/main/Speed%20Benchmark.ipynb) with StableDiffusion-XL) and uses ~45% less memory.
Here's the overview from the [project page](https://vislearn.github.io/ControlNet-XS/):
*With increasing computing capabilities, current model architectures appear to follow the trend of simply upscaling all components without validating the necessity for doing so. In this project we investigate the size and architectural design of ControlNet [Zhang et al., 2023] for controlling the image generation process with stable diffusion-based models. We show that a new architecture with as little as 1% of the parameters of the base model achieves state-of-the art results, considerably better than ControlNet in terms of FID score. Hence we call it ControlNet-XS. We provide the code for controlling StableDiffusion-XL [Podell et al., 2023] (Model B, 48M Parameters) and StableDiffusion 2.1 [Rombach et al. 2022] (Model B, 14M Parameters), all under openrail license.*
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## StableDiffusionControlNetXSPipeline
[[autodoc]] StableDiffusionControlNetXSPipeline
- all
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
@@ -1,45 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS with Stable Diffusion XL
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
ControlNet-XS generates images with comparable quality to a regular ControlNet, but it is 20-25% faster ([see benchmark](https://github.com/UmerHA/controlnet-xs-benchmark/blob/main/Speed%20Benchmark.ipynb)) and uses ~45% less memory.
Here's the overview from the [project page](https://vislearn.github.io/ControlNet-XS/):
*With increasing computing capabilities, current model architectures appear to follow the trend of simply upscaling all components without validating the necessity for doing so. In this project we investigate the size and architectural design of ControlNet [Zhang et al., 2023] for controlling the image generation process with stable diffusion-based models. We show that a new architecture with as little as 1% of the parameters of the base model achieves state-of-the art results, considerably better than ControlNet in terms of FID score. Hence we call it ControlNet-XS. We provide the code for controlling StableDiffusion-XL [Podell et al., 2023] (Model B, 48M Parameters) and StableDiffusion 2.1 [Rombach et al. 2022] (Model B, 14M Parameters), all under openrail license.*
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
<Tip warning={true}>
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) and leave us feedback on how we can improve!
</Tip>
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## StableDiffusionXLControlNetXSPipeline
[[autodoc]] StableDiffusionXLControlNetXSPipeline
- all
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
-3
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@@ -40,8 +40,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Consistency Models](consistency_models) | unconditional image generation |
| [ControlNet](controlnet) | text2image, image2image, inpainting |
| [ControlNet with Stable Diffusion XL](controlnet_sdxl) | text2image |
| [ControlNet-XS](controlnetxs) | text2image |
| [ControlNet-XS with Stable Diffusion XL](controlnetxs_sdxl) | text2image |
| [Cycle Diffusion](cycle_diffusion) | image2image |
| [Dance Diffusion](dance_diffusion) | unconditional audio generation |
| [DDIM](ddim) | unconditional image generation |
@@ -73,7 +71,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Stable Diffusion](stable_diffusion/overview) | text2image, image2image, depth2image, inpainting, image variation, latent upscaler, super-resolution |
| [Stable Diffusion Model Editing](model_editing) | model editing |
| [Stable Diffusion XL](stable_diffusion/stable_diffusion_xl) | text2image, image2image, inpainting |
| [Stable Diffusion XL Turbo](stable_diffusion/sdxl_turbo) | text2image, image2image, inpainting |
| [Stable unCLIP](stable_unclip) | text2image, image variation |
| [Stochastic Karras VE](stochastic_karras_ve) | unconditional image generation |
| [T2I-Adapter](stable_diffusion/adapter) | text2image |
@@ -20,7 +20,7 @@ The abstract from the paper is:
## Tips
- SDXL Turbo uses the exact same architecture as [SDXL](./stable_diffusion_xl), which means it also has the same API. Please refer to the [SDXL](./stable_diffusion_xl) API reference for more details.
- SDXL Turbo uses the exact same architecture as [SDXL](./stable_diffusion_xl).
- SDXL Turbo should disable guidance scale by setting `guidance_scale=0.0`
- SDXL Turbo should use `timestep_spacing='trailing'` for the scheduler and use between 1 and 4 steps.
- SDXL Turbo has been trained to generate images of size 512x512.
@@ -28,8 +28,26 @@ The abstract from the paper is:
<Tip>
To learn how to use SDXL Turbo for various tasks, how to optimize performance, and other usage examples, take a look at the [SDXL Turbo](../../../using-diffusers/sdxl_turbo) guide.
To learn how to use SDXL Turbo for various tasks, how to optimize performance, and other usage examples, take a look at the [Stable Diffusion XL](../../../using-diffusers/sdxl_turbo) guide.
Check out the [Stability AI](https://huggingface.co/stabilityai) Hub organization for the official base and refiner model checkpoints!
</Tip>
## StableDiffusionXLPipeline
[[autodoc]] StableDiffusionXLPipeline
- all
- __call__
## StableDiffusionXLImg2ImgPipeline
[[autodoc]] StableDiffusionXLImg2ImgPipeline
- all
- __call__
## StableDiffusionXLInpaintPipeline
[[autodoc]] StableDiffusionXLInpaintPipeline
- all
- __call__
@@ -485,69 +485,6 @@ image.save("sdxl_t2i.png")
</div>
</div>
You can use the IP-Adapter face model to apply specific faces to your images. It is an effective way to maintain consistent characters in your image generations.
Weights are loaded with the same method used for the other IP-Adapters.
```python
# Load ip-adapter-full-face_sd15.bin
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter-full-face_sd15.bin")
```
<Tip>
It is recommended to use `DDIMScheduler` and `EulerDiscreteScheduler` for face model.
</Tip>
```python
import torch
from diffusers import StableDiffusionPipeline, DDIMScheduler
from diffusers.utils import load_image
noise_scheduler = DDIMScheduler(
num_train_timesteps=1000,
beta_start=0.00085,
beta_end=0.012,
beta_schedule="scaled_linear",
clip_sample=False,
set_alpha_to_one=False,
steps_offset=1
)
pipeline = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
scheduler=noise_scheduler,
).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter-full-face_sd15.bin")
pipeline.set_ip_adapter_scale(0.7)
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png")
generator = torch.Generator(device="cpu").manual_seed(33)
image = pipeline(
prompt="A photo of a girl wearing a black dress, holding red roses in hand, upper body, behind is the Eiffel Tower",
ip_adapter_image=image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50, num_images_per_prompt=1, width=512, height=704,
generator=generator,
).images[0]
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ipadapter_full_face_output.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">output image</figcaption>
</div>
</div>
### LCM-Lora
@@ -133,10 +133,10 @@ def save_model_card(
diffusers_imports_pivotal = """from huggingface_hub import hf_hub_download
from safetensors.torch import load_file
"""
diffusers_example_pivotal = f"""embedding_path = hf_hub_download(repo_id='{repo_id}', filename="embeddings.safetensors", repo_type="model")
diffusers_example_pivotal = f"""embedding_path = hf_hub_download(repo_id="{repo_id}", filename="embeddings.safetensors", repo_type="model")
state_dict = load_file(embedding_path)
pipeline.load_textual_inversion(state_dict["clip_l"], token=["<s0>", "<s1>"], text_encoder=pipeline.text_encoder, tokenizer=pipeline.tokenizer)
pipeline.load_textual_inversion(state_dict["clip_g"], token=["<s0>", "<s1>"], text_encoder=pipeline.text_encoder_2, tokenizer=pipeline.tokenizer_2)
pipeline.load_textual_inversion(state_dict["clip_l"], token=["<s0>", "<s1>"], text_encoder=pipe.text_encoder, tokenizer=pipe.tokenizer)
pipeline.load_textual_inversion(state_dict["clip_g"], token=["<s0>", "<s1>"], text_encoder=pipe.text_encoder_2, tokenizer=pipe.tokenizer_2)
"""
if token_abstraction_dict:
for key, value in token_abstraction_dict.items():
@@ -145,7 +145,8 @@ pipeline.load_textual_inversion(state_dict["clip_g"], token=["<s0>", "<s1>"], te
to trigger concept `{key}` → use `{tokens}` in your prompt \n
"""
yaml = f"""---
yaml = f"""
---
tags:
- stable-diffusion-xl
- stable-diffusion-xl-diffusers
@@ -157,10 +158,8 @@ tags:
base_model: {base_model}
instance_prompt: {instance_prompt}
license: openrail++
widget:
- text: '{validation_prompt if validation_prompt else instance_prompt}'
---
"""
"""
model_card = f"""
# SDXL LoRA DreamBooth - {repo_id}
@@ -171,6 +170,14 @@ widget:
### These are {repo_id} LoRA adaption weights for {base_model}.
The weights were trained using [DreamBooth](https://dreambooth.github.io/).
LoRA for the text encoder was enabled: {train_text_encoder}.
Pivotal tuning was enabled: {train_text_encoder_ti}.
Special VAE used for training: {vae_path}.
## Trigger words
{trigger_str}
@@ -189,24 +196,11 @@ image = pipeline('{validation_prompt if validation_prompt else instance_prompt}'
For more details, including weighting, merging and fusing LoRAs, check the [documentation on loading LoRAs in diffusers](https://huggingface.co/docs/diffusers/main/en/using-diffusers/loading_adapters)
## Download model
## Download model (use it with UIs such as AUTO1111, Comfy, SD.Next, Invoke)
### Use it with UIs such as AUTOMATIC1111, Comfy UI, SD.Next, Invoke
Weights for this model are available in Safetensors format.
- Download the LoRA *.safetensors [here](/{repo_id}/blob/main/pytorch_lora_weights.safetensors). Rename it and place it on your Lora folder.
- Download the text embeddings *.safetensors [here](/{repo_id}/blob/main/embeddings.safetensors). Rename it and place it on it on your embeddings folder.
All [Files & versions](/{repo_id}/tree/main).
## Details
The weights were trained using [🧨 diffusers Advanced Dreambooth Training Script](https://github.com/huggingface/diffusers/blob/main/examples/advanced_diffusion_training/train_dreambooth_lora_sdxl_advanced.py).
LoRA for the text encoder was enabled. {train_text_encoder}.
Pivotal tuning was enabled: {train_text_encoder_ti}.
Special VAE used for training: {vae_path}.
[Download]({repo_id}/tree/main) them in the Files & versions tab.
"""
with open(os.path.join(repo_folder, "README.md"), "w") as f:
@@ -673,12 +667,6 @@ def parse_args(input_args=None):
default=4,
help=("The dimension of the LoRA update matrices."),
)
parser.add_argument(
"--cache_latents",
action="store_true",
default=False,
help="Cache the VAE latents",
)
if input_args is not None:
args = parser.parse_args(input_args)
@@ -1182,7 +1170,6 @@ def main(args):
revision=args.revision,
variant=args.variant,
)
vae_scaling_factor = vae.config.scaling_factor
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
@@ -1613,20 +1600,6 @@ def main(args):
args.validation_prompt = args.validation_prompt.replace(token_abs, "".join(token_replacement))
print("validation prompt:", args.validation_prompt)
if args.cache_latents:
latents_cache = []
for batch in tqdm(train_dataloader, desc="Caching latents"):
with torch.no_grad():
batch["pixel_values"] = batch["pixel_values"].to(
accelerator.device, non_blocking=True, dtype=torch.float32
)
latents_cache.append(vae.encode(batch["pixel_values"]).latent_dist)
if args.validation_prompt is None:
del vae
if torch.cuda.is_available():
torch.cuda.empty_cache()
# Scheduler and math around the number of training steps.
overrode_max_train_steps = False
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
@@ -1742,7 +1715,9 @@ def main(args):
unet.train()
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
pixel_values = batch["pixel_values"].to(dtype=vae.dtype)
prompts = batch["prompts"]
# print(prompts)
# encode batch prompts when custom prompts are provided for each image -
if train_dataset.custom_instance_prompts:
if freeze_text_encoder:
@@ -1754,13 +1729,9 @@ def main(args):
tokens_one = tokenize_prompt(tokenizer_one, prompts, add_special_tokens)
tokens_two = tokenize_prompt(tokenizer_two, prompts, add_special_tokens)
if args.cache_latents:
model_input = latents_cache[step].sample()
else:
pixel_values = batch["pixel_values"].to(dtype=vae.dtype)
model_input = vae.encode(pixel_values).latent_dist.sample()
model_input = model_input * vae_scaling_factor
# Convert images to latent space
model_input = vae.encode(pixel_values).latent_dist.sample()
model_input = model_input * vae.config.scaling_factor
if args.pretrained_vae_model_name_or_path is None:
model_input = model_input.to(weight_dtype)
-65
View File
@@ -50,7 +50,6 @@ prompt-to-prompt | change parts of a prompt and retain image structure (see [pap
| Latent Consistency Interpolation Pipeline | Interpolate the latent space of Latent Consistency Models with multiple prompts | [Latent Consistency Interpolation Pipeline](#latent-consistency-interpolation-pipeline) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1pK3NrLWJSiJsBynLns1K1-IDTW9zbPvl?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) |
| Regional Prompting Pipeline | Assign multiple prompts for different regions | [Regional Prompting Pipeline](#regional-prompting-pipeline) | - | [hako-mikan](https://github.com/hako-mikan) |
| LDM3D-sr (LDM3D upscaler) | Upscale low resolution RGB and depth inputs to high resolution | [StableDiffusionUpscaleLDM3D Pipeline](https://github.com/estelleafl/diffusers/tree/ldm3d_upscaler_community/examples/community#stablediffusionupscaleldm3d-pipeline) | - | [Estelle Aflalo](https://github.com/estelleafl) |
| AnimateDiff ControlNet Pipeline | Combines AnimateDiff with precise motion control using ControlNets | [AnimateDiff ControlNet Pipeline](#animatediff-controlnet-pipeline) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1SKboYeGjEQmQPWoFC0aLYpBlYdHXkvAu?usp=sharing) | [Aryan V S](https://github.com/a-r-r-o-w) and [Edoardo Botta](https://github.com/EdoardoBotta) |
| DemoFusion Pipeline | Implementation of [DemoFusion: Democratising High-Resolution Image Generation With No $$$](https://arxiv.org/abs/2311.16973) | [DemoFusion Pipeline](#DemoFusion) | - | [Ruoyi Du](https://github.com/RuoyiDu) |
To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly.
@@ -2840,70 +2839,6 @@ The Pipeline supports `compel` syntax. Input prompts using the `compel` structur
* Reconstructed image:
* ![dps_generated_image](https://github.com/tongdaxu/Images/assets/22267548/b74f084d-93f4-4845-83d8-44c0fa758a5f)
### AnimateDiff ControlNet Pipeline
This pipeline combines AnimateDiff and ControlNet. Enjoy precise motion control for your videos! Refer to [this](https://github.com/huggingface/diffusers/issues/5866) issue for more details.
```py
import torch
from diffusers import AutoencoderKL, ControlNetModel, MotionAdapter
from diffusers.pipelines import DiffusionPipeline
from diffusers.schedulers import DPMSolverMultistepScheduler
from PIL import Image
motion_id = "guoyww/animatediff-motion-adapter-v1-5-2"
adapter = MotionAdapter.from_pretrained(motion_id)
controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_openpose", torch_dtype=torch.float16)
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16)
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
pipe = DiffusionPipeline.from_pretrained(
model_id,
motion_adapter=adapter,
controlnet=controlnet,
vae=vae,
custom_pipeline="pipeline_animatediff_controlnet",
).to(device="cuda", dtype=torch.float16)
pipe.scheduler = DPMSolverMultistepScheduler.from_pretrained(
model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", steps_offset=1
)
pipe.enable_vae_slicing()
conditioning_frames = []
for i in range(1, 16 + 1):
conditioning_frames.append(Image.open(f"frame_{i}.png"))
prompt = "astronaut in space, dancing"
negative_prompt = "bad quality, worst quality, jpeg artifacts, ugly"
result = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
width=512,
height=768,
conditioning_frames=conditioning_frames,
num_inference_steps=12,
).frames[0]
from diffusers.utils import export_to_gif
export_to_gif(result.frames[0], "result.gif")
```
<table>
<tr><td colspan="2" align=center><b>Conditioning Frames</b></td></tr>
<tr align=center>
<td align=center><img src="https://user-images.githubusercontent.com/7365912/265043418-23291941-864d-495a-8ba8-d02e05756396.gif" alt="input-frames"></td>
</tr>
<tr><td colspan="2" align=center><b>AnimateDiff model: SG161222/Realistic_Vision_V5.1_noVAE</b></td></tr>
<tr>
<td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/baf301e2-d03c-4129-bd84-203a1de2b2be" alt="gif-1"></td>
<td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/9f923475-ecaf-452b-92c8-4e42171182d8" alt="gif-2"></td>
</tr>
<tr><td colspan="2" align=center><b>AnimateDiff model: CardosAnime</b></td></tr>
<tr>
<td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/b2c41028-38a0-45d6-86ed-fec7446b87f7" alt="gif-1"></td>
<td align=center><img src="https://github.com/huggingface/diffusers/assets/72266394/eb7d2952-72e4-44fa-b664-077c79b4fc70" alt="gif-2"></td>
</tr>
</table>
### DemoFusion
This pipeline is the official implementation of [DemoFusion: Democratising High-Resolution Image Generation With No $$$](https://arxiv.org/abs/2311.16973).
The original repo can be found at [repo](https://github.com/PRIS-CV/DemoFusion).
File diff suppressed because it is too large Load Diff
+14 -22
View File
@@ -1,6 +1,6 @@
# Latent Consistency Distillation Example:
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is a method to distill a latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use latent consistency distillation to distill stable-diffusion-v1.5 for inference with few timesteps.
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is method to distill latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use the latent consistency distillation to distill stable-diffusion-v1.5 for less timestep inference.
## Full model distillation
@@ -24,7 +24,7 @@ Then cd in the example folder and run
pip install -r requirements.txt
```
And initialize an [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
@@ -46,16 +46,12 @@ write_basic_config()
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
#### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example, and for illustrative purposes only. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/). You may also need to search the hyperparameter space according to the dataset you use.
#### Example with LAION-A6+ dataset
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_sd_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
runwayml/stable-diffusion-v1-5
PROGRAM="train_lcm_distill_sd_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
--resolution=512 \
@@ -63,7 +59,7 @@ accelerate launch train_lcm_distill_sd_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -73,23 +69,19 @@ accelerate launch train_lcm_distill_sd_wds.py \
--resume_from_checkpoint=latest \
--report_to=wandb \
--seed=453645634 \
--push_to_hub
--push_to_hub \
```
## LCM-LoRA
Instead of fine-tuning the full model, we can also just train a LoRA that can be injected into any SDXL model.
### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/).
### Example with LAION-A6+ dataset
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_lora_sd_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
runwayml/stable-diffusion-v1-5
PROGRAM="train_lcm_distill_lora_sd_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
--resolution=512 \
@@ -98,7 +90,7 @@ accelerate launch train_lcm_distill_lora_sd_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -1,6 +1,6 @@
# Latent Consistency Distillation Example:
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is a method to distill a latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use latent consistency distillation to distill SDXL for inference with few timesteps.
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is method to distill latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use the latent consistency distillation to distill SDXL for less timestep inference.
## Full model distillation
@@ -24,7 +24,7 @@ Then cd in the example folder and run
pip install -r requirements.txt
```
And initialize an [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
@@ -46,16 +46,12 @@ write_basic_config()
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
#### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example, and for illustrative purposes only. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/). You may also need to search the hyperparameter space according to the dataset you use.
#### Example with LAION-A6+ dataset
```bash
export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
export MODEL_DIR="stabilityai/stable-diffusion-xl-base-1.0"
PROGRAM="train_lcm_distill_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--pretrained_vae_model_name_or_path=madebyollin/sdxl-vae-fp16-fix \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
@@ -64,7 +60,7 @@ accelerate launch train_lcm_distill_sdxl_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -81,15 +77,11 @@ accelerate launch train_lcm_distill_sdxl_wds.py \
Instead of fine-tuning the full model, we can also just train a LoRA that can be injected into any SDXL model.
### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/).
### Example with LAION-A6+ dataset
```bash
export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_lora_sdxl_wds.py \
export MODEL_DIR="stabilityai/stable-diffusion-xl-base-1.0"
PROGRAM="train_lcm_distill_lora_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--pretrained_vae_model_name_or_path=madebyollin/sdxl-vae-fp16-fix \
--output_dir=$OUTPUT_DIR \
@@ -100,7 +92,7 @@ accelerate launch train_lcm_distill_lora_sdxl_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -1123,7 +1123,7 @@ def main(args):
for epoch in range(first_epoch, args.num_train_epochs):
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
image, text = batch
image, text, _, _ = batch
image = image.to(accelerator.device, non_blocking=True)
encoded_text = compute_embeddings_fn(text)
@@ -68,11 +68,6 @@ from diffusers.utils.import_utils import is_xformers_available
MAX_SEQ_LENGTH = 77
# Adjust for your dataset
WDS_JSON_WIDTH = "width" # original_width for LAION
WDS_JSON_HEIGHT = "height" # original_height for LAION
MIN_SIZE = 700 # ~960 for LAION, ideal: 1024 if the dataset contains large images
if is_wandb_available():
import wandb
@@ -151,10 +146,10 @@ class WebdatasetFilter:
try:
if "json" in x:
x_json = json.loads(x["json"])
filter_size = (x_json.get(WDS_JSON_WIDTH, 0.0) or 0.0) >= self.min_size and x_json.get(
WDS_JSON_HEIGHT, 0
filter_size = (x_json.get("original_width", 0.0) or 0.0) >= self.min_size and x_json.get(
"original_height", 0
) >= self.min_size
filter_watermark = (x_json.get("pwatermark", 0.0) or 0.0) <= self.max_pwatermark
filter_watermark = (x_json.get("pwatermark", 1.0) or 1.0) <= self.max_pwatermark
return filter_size and filter_watermark
else:
return False
@@ -185,7 +180,7 @@ class Text2ImageDataset:
if use_fix_crop_and_size:
return (resolution, resolution)
else:
return (int(json.get(WDS_JSON_WIDTH, 0.0)), int(json.get(WDS_JSON_HEIGHT, 0.0)))
return (int(json.get("original_width", 0.0)), int(json.get("original_height", 0.0)))
def transform(example):
# resize image
@@ -217,7 +212,7 @@ class Text2ImageDataset:
pipeline = [
wds.ResampledShards(train_shards_path_or_url),
tarfile_to_samples_nothrow,
wds.select(WebdatasetFilter(min_size=MIN_SIZE)),
wds.select(WebdatasetFilter(min_size=960)),
wds.shuffle(shuffle_buffer_size),
*processing_pipeline,
wds.batched(per_gpu_batch_size, partial=False, collation_fn=default_collate),
@@ -1106,7 +1106,7 @@ def main(args):
for epoch in range(first_epoch, args.num_train_epochs):
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
image, text = batch
image, text, _, _ = batch
image = image.to(accelerator.device, non_blocking=True)
encoded_text = compute_embeddings_fn(text)
@@ -67,11 +67,6 @@ from diffusers.utils.import_utils import is_xformers_available
MAX_SEQ_LENGTH = 77
# Adjust for your dataset
WDS_JSON_WIDTH = "width" # original_width for LAION
WDS_JSON_HEIGHT = "height" # original_height for LAION
MIN_SIZE = 700 # ~960 for LAION, ideal: 1024 if the dataset contains large images
if is_wandb_available():
import wandb
@@ -133,10 +128,10 @@ class WebdatasetFilter:
try:
if "json" in x:
x_json = json.loads(x["json"])
filter_size = (x_json.get(WDS_JSON_WIDTH, 0.0) or 0.0) >= self.min_size and x_json.get(
WDS_JSON_HEIGHT, 0
filter_size = (x_json.get("original_width", 0.0) or 0.0) >= self.min_size and x_json.get(
"original_height", 0
) >= self.min_size
filter_watermark = (x_json.get("pwatermark", 0.0) or 0.0) <= self.max_pwatermark
filter_watermark = (x_json.get("pwatermark", 1.0) or 1.0) <= self.max_pwatermark
return filter_size and filter_watermark
else:
return False
@@ -167,7 +162,7 @@ class Text2ImageDataset:
if use_fix_crop_and_size:
return (resolution, resolution)
else:
return (int(json.get(WDS_JSON_WIDTH, 0.0)), int(json.get(WDS_JSON_HEIGHT, 0.0)))
return (int(json.get("original_width", 0.0)), int(json.get("original_height", 0.0)))
def transform(example):
# resize image
@@ -199,7 +194,7 @@ class Text2ImageDataset:
pipeline = [
wds.ResampledShards(train_shards_path_or_url),
tarfile_to_samples_nothrow,
wds.select(WebdatasetFilter(min_size=MIN_SIZE)),
wds.select(WebdatasetFilter(min_size=960)),
wds.shuffle(shuffle_buffer_size),
*processing_pipeline,
wds.batched(per_gpu_batch_size, partial=False, collation_fn=default_collate),
-1
View File
@@ -44,7 +44,6 @@ write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Dog toy example
-1
View File
@@ -47,7 +47,6 @@ write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Dog toy example
-1
View File
@@ -4,4 +4,3 @@ transformers>=4.25.1
ftfy
tensorboard
Jinja2
peft==0.7.0
@@ -4,4 +4,3 @@ transformers>=4.25.1
ftfy
tensorboard
Jinja2
peft==0.7.0
+99 -27
View File
@@ -16,6 +16,7 @@
import argparse
import copy
import gc
import itertools
import logging
import math
import os
@@ -34,8 +35,6 @@ from accelerate.utils import ProjectConfiguration, set_seed
from huggingface_hub import create_repo, upload_folder
from huggingface_hub.utils import insecure_hashlib
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from PIL import Image
from PIL.ImageOps import exif_transpose
from torch.utils.data import Dataset
@@ -53,7 +52,14 @@ from diffusers import (
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.models.attention_processor import (
AttnAddedKVProcessor,
AttnAddedKVProcessor2_0,
SlicedAttnAddedKVProcessor,
)
from diffusers.models.lora import LoRALinearLayer
from diffusers.optimization import get_scheduler
from diffusers.training_utils import unet_lora_state_dict
from diffusers.utils import check_min_version, is_wandb_available
from diffusers.utils.import_utils import is_xformers_available
@@ -858,19 +864,79 @@ def main(args):
text_encoder.gradient_checkpointing_enable()
# now we will add new LoRA weights to the attention layers
unet_lora_config = LoraConfig(
r=args.rank,
init_lora_weights="gaussian",
target_modules=["to_k", "to_q", "to_v", "to_out.0", "add_k_proj", "add_v_proj"],
)
unet.add_adapter(unet_lora_config)
# It's important to realize here how many attention weights will be added and of which sizes
# The sizes of the attention layers consist only of two different variables:
# 1) - the "hidden_size", which is increased according to `unet.config.block_out_channels`.
# 2) - the "cross attention size", which is set to `unet.config.cross_attention_dim`.
# The text encoder comes from 🤗 transformers, we will also attach adapters to it.
if args.train_text_encoder:
text_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["q_proj", "k_proj", "v_proj", "out_proj"]
# Let's first see how many attention processors we will have to set.
# For Stable Diffusion, it should be equal to:
# - down blocks (2x attention layers) * (2x transformer layers) * (3x down blocks) = 12
# - mid blocks (2x attention layers) * (1x transformer layers) * (1x mid blocks) = 2
# - up blocks (2x attention layers) * (3x transformer layers) * (3x up blocks) = 18
# => 32 layers
# Set correct lora layers
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
text_encoder.add_adapter(text_lora_config)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
if isinstance(attn_processor, (AttnAddedKVProcessor, SlicedAttnAddedKVProcessor, AttnAddedKVProcessor2_0)):
attn_module.add_k_proj.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.add_k_proj.in_features,
out_features=attn_module.add_k_proj.out_features,
rank=args.rank,
)
)
attn_module.add_v_proj.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.add_v_proj.in_features,
out_features=attn_module.add_v_proj.out_features,
rank=args.rank,
)
)
unet_lora_parameters.extend(attn_module.add_k_proj.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.add_v_proj.lora_layer.parameters())
# The text encoder comes from 🤗 transformers, so we cannot directly modify it.
# So, instead, we monkey-patch the forward calls of its attention-blocks.
if args.train_text_encoder:
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
text_lora_parameters = LoraLoaderMixin._modify_text_encoder(text_encoder, dtype=torch.float32, rank=args.rank)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
def save_model_hook(models, weights, output_dir):
@@ -882,9 +948,9 @@ def main(args):
for model in models:
if isinstance(model, type(accelerator.unwrap_model(unet))):
unet_lora_layers_to_save = get_peft_model_state_dict(model)
unet_lora_layers_to_save = unet_lora_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder))):
text_encoder_lora_layers_to_save = get_peft_model_state_dict(model)
text_encoder_lora_layers_to_save = text_encoder_lora_state_dict(model)
else:
raise ValueError(f"unexpected save model: {model.__class__}")
@@ -944,10 +1010,11 @@ def main(args):
optimizer_class = torch.optim.AdamW
# Optimizer creation
params_to_optimize = list(filter(lambda p: p.requires_grad, unet.parameters()))
if args.train_text_encoder:
params_to_optimize = params_to_optimize + list(filter(lambda p: p.requires_grad, text_encoder.parameters()))
params_to_optimize = (
itertools.chain(unet_lora_parameters, text_lora_parameters)
if args.train_text_encoder
else unet_lora_parameters
)
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
@@ -1190,7 +1257,12 @@ def main(args):
accelerator.backward(loss)
if accelerator.sync_gradients:
accelerator.clip_grad_norm_(params_to_optimize, args.max_grad_norm)
params_to_clip = (
itertools.chain(unet_lora_parameters, text_lora_parameters)
if args.train_text_encoder
else unet_lora_parameters
)
accelerator.clip_grad_norm_(params_to_clip, args.max_grad_norm)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
@@ -1313,19 +1385,19 @@ def main(args):
if accelerator.is_main_process:
unet = accelerator.unwrap_model(unet)
unet = unet.to(torch.float32)
unet_lora_layers = unet_lora_state_dict(unet)
unet_lora_state_dict = get_peft_model_state_dict(unet)
if args.train_text_encoder:
if text_encoder is not None and args.train_text_encoder:
text_encoder = accelerator.unwrap_model(text_encoder)
text_encoder_state_dict = get_peft_model_state_dict(text_encoder)
text_encoder = text_encoder.to(torch.float32)
text_encoder_lora_layers = text_encoder_lora_state_dict(text_encoder)
else:
text_encoder_state_dict = None
text_encoder_lora_layers = None
LoraLoaderMixin.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_state_dict,
text_encoder_lora_layers=text_encoder_state_dict,
unet_lora_layers=unet_lora_layers,
text_encoder_lora_layers=text_encoder_lora_layers,
)
# Final inference
@@ -34,8 +34,6 @@ from accelerate.utils import DistributedDataParallelKwargs, ProjectConfiguration
from huggingface_hub import create_repo, upload_folder
from huggingface_hub.utils import insecure_hashlib
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from PIL import Image
from PIL.ImageOps import exif_transpose
from torch.utils.data import Dataset
@@ -52,8 +50,9 @@ from diffusers import (
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.models.lora import LoRALinearLayer
from diffusers.optimization import get_scheduler
from diffusers.training_utils import compute_snr
from diffusers.training_utils import compute_snr, unet_lora_state_dict
from diffusers.utils import check_min_version, is_wandb_available
from diffusers.utils.import_utils import is_xformers_available
@@ -1010,19 +1009,54 @@ def main(args):
text_encoder_two.gradient_checkpointing_enable()
# now we will add new LoRA weights to the attention layers
unet_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["to_k", "to_q", "to_v", "to_out.0"]
)
unet.add_adapter(unet_lora_config)
# Set correct lora layers
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
# The text encoder comes from 🤗 transformers, so we cannot directly modify it.
# So, instead, we monkey-patch the forward calls of its attention-blocks.
if args.train_text_encoder:
text_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["q_proj", "k_proj", "v_proj", "out_proj"]
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
text_lora_parameters_one = LoraLoaderMixin._modify_text_encoder(
text_encoder_one, dtype=torch.float32, rank=args.rank
)
text_lora_parameters_two = LoraLoaderMixin._modify_text_encoder(
text_encoder_two, dtype=torch.float32, rank=args.rank
)
text_encoder_one.add_adapter(text_lora_config)
text_encoder_two.add_adapter(text_lora_config)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
def save_model_hook(models, weights, output_dir):
@@ -1035,11 +1069,11 @@ def main(args):
for model in models:
if isinstance(model, type(accelerator.unwrap_model(unet))):
unet_lora_layers_to_save = get_peft_model_state_dict(model)
unet_lora_layers_to_save = unet_lora_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_one))):
text_encoder_one_lora_layers_to_save = get_peft_model_state_dict(model)
text_encoder_one_lora_layers_to_save = text_encoder_lora_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_two))):
text_encoder_two_lora_layers_to_save = get_peft_model_state_dict(model)
text_encoder_two_lora_layers_to_save = text_encoder_lora_state_dict(model)
else:
raise ValueError(f"unexpected save model: {model.__class__}")
@@ -1096,12 +1130,6 @@ def main(args):
args.learning_rate * args.gradient_accumulation_steps * args.train_batch_size * accelerator.num_processes
)
unet_lora_parameters = list(filter(lambda p: p.requires_grad, unet.parameters()))
if args.train_text_encoder:
text_lora_parameters_one = list(filter(lambda p: p.requires_grad, text_encoder_one.parameters()))
text_lora_parameters_two = list(filter(lambda p: p.requires_grad, text_encoder_two.parameters()))
# Optimization parameters
unet_lora_parameters_with_lr = {"params": unet_lora_parameters, "lr": args.learning_rate}
if args.train_text_encoder:
@@ -1166,10 +1194,26 @@ def main(args):
optimizer_class = prodigyopt.Prodigy
if args.learning_rate <= 0.1:
logger.warn(
"Learning rate is too low. When using prodigy, it's generally better to set learning rate around 1.0"
)
if args.train_text_encoder and args.text_encoder_lr:
logger.warn(
f"Learning rates were provided both for the unet and the text encoder- e.g. text_encoder_lr:"
f" {args.text_encoder_lr} and learning_rate: {args.learning_rate}. "
f"When using prodigy only learning_rate is used as the initial learning rate."
)
# changes the learning rate of text_encoder_parameters_one and text_encoder_parameters_two to be
# --learning_rate
params_to_optimize[1]["lr"] = args.learning_rate
params_to_optimize[2]["lr"] = args.learning_rate
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
beta3=args.prodigy_beta3,
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
decouple=args.prodigy_decouple,
@@ -1615,13 +1659,13 @@ def main(args):
if accelerator.is_main_process:
unet = accelerator.unwrap_model(unet)
unet = unet.to(torch.float32)
unet_lora_layers = get_peft_model_state_dict(unet)
unet_lora_layers = unet_lora_state_dict(unet)
if args.train_text_encoder:
text_encoder_one = accelerator.unwrap_model(text_encoder_one)
text_encoder_lora_layers = get_peft_model_state_dict(text_encoder_one.to(torch.float32))
text_encoder_lora_layers = text_encoder_lora_state_dict(text_encoder_one.to(torch.float32))
text_encoder_two = accelerator.unwrap_model(text_encoder_two)
text_encoder_2_lora_layers = get_peft_model_state_dict(text_encoder_two.to(torch.float32))
text_encoder_2_lora_layers = text_encoder_lora_state_dict(text_encoder_two.to(torch.float32))
else:
text_encoder_lora_layers = None
text_encoder_2_lora_layers = None
-2
View File
@@ -32,8 +32,6 @@ And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) e
accelerate config
```
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Pokemon example
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-4`, so you'll need to visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4), read the license and tick the checkbox if you agree.
-1
View File
@@ -45,7 +45,6 @@ write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Training
-1
View File
@@ -5,4 +5,3 @@ datasets
ftfy
tensorboard
Jinja2
peft==0.7.0
@@ -5,4 +5,3 @@ ftfy
tensorboard
Jinja2
datasets
peft==0.7.0
@@ -34,14 +34,13 @@ from accelerate.utils import ProjectConfiguration, set_seed
from datasets import load_dataset
from huggingface_hub import create_repo, upload_folder
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
import diffusers
from diffusers import AutoencoderKL, DDPMScheduler, DiffusionPipeline, StableDiffusionPipeline, UNet2DConditionModel
from diffusers import AutoencoderKL, DDPMScheduler, DiffusionPipeline, UNet2DConditionModel
from diffusers.models.lora import LoRALinearLayer
from diffusers.optimization import get_scheduler
from diffusers.training_utils import compute_snr
from diffusers.utils import check_min_version, is_wandb_available
@@ -480,20 +479,62 @@ def main():
elif accelerator.mixed_precision == "bf16":
weight_dtype = torch.bfloat16
# Freeze the unet parameters before adding adapters
for param in unet.parameters():
param.requires_grad_(False)
unet_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["to_k", "to_q", "to_v", "to_out.0"]
)
# Move unet, vae and text_encoder to device and cast to weight_dtype
unet.to(accelerator.device, dtype=weight_dtype)
vae.to(accelerator.device, dtype=weight_dtype)
text_encoder.to(accelerator.device, dtype=weight_dtype)
unet.add_adapter(unet_lora_config)
# now we will add new LoRA weights to the attention layers
# It's important to realize here how many attention weights will be added and of which sizes
# The sizes of the attention layers consist only of two different variables:
# 1) - the "hidden_size", which is increased according to `unet.config.block_out_channels`.
# 2) - the "cross attention size", which is set to `unet.config.cross_attention_dim`.
# Let's first see how many attention processors we will have to set.
# For Stable Diffusion, it should be equal to:
# - down blocks (2x attention layers) * (2x transformer layers) * (3x down blocks) = 12
# - mid blocks (2x attention layers) * (1x transformer layers) * (1x mid blocks) = 2
# - up blocks (2x attention layers) * (3x transformer layers) * (3x down blocks) = 18
# => 32 layers
# Set correct lora layers
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
if args.enable_xformers_memory_efficient_attention:
if is_xformers_available():
@@ -508,8 +549,6 @@ def main():
else:
raise ValueError("xformers is not available. Make sure it is installed correctly")
lora_layers = filter(lambda p: p.requires_grad, unet.parameters())
# Enable TF32 for faster training on Ampere GPUs,
# cf https://pytorch.org/docs/stable/notes/cuda.html#tensorfloat-32-tf32-on-ampere-devices
if args.allow_tf32:
@@ -534,7 +573,7 @@ def main():
optimizer_cls = torch.optim.AdamW
optimizer = optimizer_cls(
lora_layers,
unet_lora_parameters,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
@@ -661,8 +700,8 @@ def main():
)
# Prepare everything with our `accelerator`.
unet, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet, optimizer, train_dataloader, lr_scheduler
unet_lora_parameters, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet_lora_parameters, optimizer, train_dataloader, lr_scheduler
)
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
@@ -794,7 +833,7 @@ def main():
# Backpropagate
accelerator.backward(loss)
if accelerator.sync_gradients:
params_to_clip = lora_layers
params_to_clip = unet_lora_parameters
accelerator.clip_grad_norm_(params_to_clip, args.max_grad_norm)
optimizer.step()
lr_scheduler.step()
@@ -831,15 +870,6 @@ def main():
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
accelerator.save_state(save_path)
unet_lora_state_dict = get_peft_model_state_dict(unet)
StableDiffusionPipeline.save_lora_weights(
save_directory=save_path,
unet_lora_layers=unet_lora_state_dict,
safe_serialization=True,
)
logger.info(f"Saved state to {save_path}")
logs = {"step_loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
@@ -896,13 +926,7 @@ def main():
accelerator.wait_for_everyone()
if accelerator.is_main_process:
unet = unet.to(torch.float32)
unet_lora_state_dict = get_peft_model_state_dict(unet)
StableDiffusionPipeline.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_state_dict,
safe_serialization=True,
)
unet.save_attn_procs(args.output_dir)
if args.push_to_hub:
save_model_card(
@@ -16,6 +16,7 @@
"""Fine-tuning script for Stable Diffusion XL for text2image with support for LoRA."""
import argparse
import itertools
import logging
import math
import os
@@ -36,8 +37,6 @@ from accelerate.utils import DistributedDataParallelKwargs, ProjectConfiguration
from datasets import load_dataset
from huggingface_hub import create_repo, upload_folder
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from torchvision import transforms
from torchvision.transforms.functional import crop
from tqdm.auto import tqdm
@@ -51,6 +50,7 @@ from diffusers import (
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.models.lora import LoRALinearLayer
from diffusers.optimization import get_scheduler
from diffusers.training_utils import compute_snr
from diffusers.utils import check_min_version, is_wandb_available
@@ -658,20 +658,53 @@ def main(args):
# now we will add new LoRA weights to the attention layers
# Set correct lora layers
unet_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["to_k", "to_q", "to_v", "to_out.0"]
)
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
unet.add_adapter(unet_lora_config)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# The text encoder comes from 🤗 transformers, we will also attach adapters to it.
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
# The text encoder comes from 🤗 transformers, so we cannot directly modify it.
# So, instead, we monkey-patch the forward calls of its attention-blocks.
if args.train_text_encoder:
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
text_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["q_proj", "k_proj", "v_proj", "out_proj"]
text_lora_parameters_one = LoraLoaderMixin._modify_text_encoder(
text_encoder_one, dtype=torch.float32, rank=args.rank
)
text_lora_parameters_two = LoraLoaderMixin._modify_text_encoder(
text_encoder_two, dtype=torch.float32, rank=args.rank
)
text_encoder_one.add_adapter(text_lora_config)
text_encoder_two.add_adapter(text_lora_config)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
def save_model_hook(models, weights, output_dir):
@@ -684,11 +717,11 @@ def main(args):
for model in models:
if isinstance(model, type(accelerator.unwrap_model(unet))):
unet_lora_layers_to_save = get_peft_model_state_dict(model)
unet_lora_layers_to_save = unet_attn_processors_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_one))):
text_encoder_one_lora_layers_to_save = get_peft_model_state_dict(model)
text_encoder_one_lora_layers_to_save = text_encoder_lora_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_two))):
text_encoder_two_lora_layers_to_save = get_peft_model_state_dict(model)
text_encoder_two_lora_layers_to_save = text_encoder_lora_state_dict(model)
else:
raise ValueError(f"unexpected save model: {model.__class__}")
@@ -759,13 +792,11 @@ def main(args):
optimizer_class = torch.optim.AdamW
# Optimizer creation
params_to_optimize = list(filter(lambda p: p.requires_grad, unet.parameters()))
if args.train_text_encoder:
params_to_optimize = (
params_to_optimize
+ list(filter(lambda p: p.requires_grad, text_encoder_one.parameters()))
+ list(filter(lambda p: p.requires_grad, text_encoder_two.parameters()))
)
params_to_optimize = (
itertools.chain(unet_lora_parameters, text_lora_parameters_one, text_lora_parameters_two)
if args.train_text_encoder
else unet_lora_parameters
)
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
@@ -1097,7 +1128,12 @@ def main(args):
# Backpropagate
accelerator.backward(loss)
if accelerator.sync_gradients:
accelerator.clip_grad_norm_(params_to_optimize, args.max_grad_norm)
params_to_clip = (
itertools.chain(unet_lora_parameters, text_lora_parameters_one, text_lora_parameters_two)
if args.train_text_encoder
else unet_lora_parameters
)
accelerator.clip_grad_norm_(params_to_clip, args.max_grad_norm)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
@@ -1193,21 +1229,20 @@ def main(args):
accelerator.wait_for_everyone()
if accelerator.is_main_process:
unet = accelerator.unwrap_model(unet)
unet_lora_state_dict = get_peft_model_state_dict(unet)
unet_lora_layers = unet_attn_processors_state_dict(unet)
if args.train_text_encoder:
text_encoder_one = accelerator.unwrap_model(text_encoder_one)
text_encoder_lora_layers = text_encoder_lora_state_dict(text_encoder_one)
text_encoder_two = accelerator.unwrap_model(text_encoder_two)
text_encoder_lora_layers = get_peft_model_state_dict(text_encoder_one)
text_encoder_2_lora_layers = get_peft_model_state_dict(text_encoder_two)
text_encoder_2_lora_layers = text_encoder_lora_state_dict(text_encoder_two)
else:
text_encoder_lora_layers = None
text_encoder_2_lora_layers = None
StableDiffusionXLPipeline.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_state_dict,
unet_lora_layers=unet_lora_layers,
text_encoder_lora_layers=text_encoder_lora_layers,
text_encoder_2_lora_layers=text_encoder_2_lora_layers,
)
-6
View File
@@ -80,7 +80,6 @@ else:
"AutoencoderTiny",
"ConsistencyDecoderVAE",
"ControlNetModel",
"ControlNetXSModel",
"Kandinsky3UNet",
"ModelMixin",
"MotionAdapter",
@@ -251,7 +250,6 @@ else:
"StableDiffusionControlNetImg2ImgPipeline",
"StableDiffusionControlNetInpaintPipeline",
"StableDiffusionControlNetPipeline",
"StableDiffusionControlNetXSPipeline",
"StableDiffusionDepth2ImgPipeline",
"StableDiffusionDiffEditPipeline",
"StableDiffusionGLIGENPipeline",
@@ -275,7 +273,6 @@ else:
"StableDiffusionXLControlNetImg2ImgPipeline",
"StableDiffusionXLControlNetInpaintPipeline",
"StableDiffusionXLControlNetPipeline",
"StableDiffusionXLControlNetXSPipeline",
"StableDiffusionXLImg2ImgPipeline",
"StableDiffusionXLInpaintPipeline",
"StableDiffusionXLInstructPix2PixPipeline",
@@ -457,7 +454,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
AutoencoderTiny,
ConsistencyDecoderVAE,
ControlNetModel,
ControlNetXSModel,
Kandinsky3UNet,
ModelMixin,
MotionAdapter,
@@ -607,7 +603,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
StableDiffusionControlNetImg2ImgPipeline,
StableDiffusionControlNetInpaintPipeline,
StableDiffusionControlNetPipeline,
StableDiffusionControlNetXSPipeline,
StableDiffusionDepth2ImgPipeline,
StableDiffusionDiffEditPipeline,
StableDiffusionGLIGENPipeline,
@@ -631,7 +626,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
StableDiffusionXLControlNetImg2ImgPipeline,
StableDiffusionXLControlNetInpaintPipeline,
StableDiffusionXLControlNetPipeline,
StableDiffusionXLControlNetXSPipeline,
StableDiffusionXLImg2ImgPipeline,
StableDiffusionXLInpaintPipeline,
StableDiffusionXLInstructPix2PixPipeline,
+2 -29
View File
@@ -22,7 +22,7 @@ import torch.nn.functional as F
from huggingface_hub.utils import validate_hf_hub_args
from torch import nn
from ..models.embeddings import ImageProjection, MLPProjection, Resampler
from ..models.embeddings import ImageProjection, Resampler
from ..models.modeling_utils import _LOW_CPU_MEM_USAGE_DEFAULT, load_model_dict_into_meta
from ..utils import (
USE_PEFT_BACKEND,
@@ -675,9 +675,6 @@ class UNet2DConditionLoadersMixin:
if "proj.weight" in state_dict["image_proj"]:
# IP-Adapter
num_image_text_embeds = 4
elif "proj.3.weight" in state_dict["image_proj"]:
# IP-Adapter Full Face
num_image_text_embeds = 257 # 256 CLIP tokens + 1 CLS token
else:
# IP-Adapter Plus
num_image_text_embeds = state_dict["image_proj"]["latents"].shape[1]
@@ -747,32 +744,8 @@ class UNet2DConditionLoadersMixin:
"norm.bias": state_dict["image_proj"]["norm.bias"],
}
)
image_projection.load_state_dict(image_proj_state_dict)
del image_proj_state_dict
elif "proj.3.weight" in state_dict["image_proj"]:
clip_embeddings_dim = state_dict["image_proj"]["proj.0.weight"].shape[0]
cross_attention_dim = state_dict["image_proj"]["proj.3.weight"].shape[0]
image_projection = MLPProjection(
cross_attention_dim=cross_attention_dim, image_embed_dim=clip_embeddings_dim
)
image_projection.to(dtype=self.dtype, device=self.device)
# load image projection layer weights
image_proj_state_dict = {}
image_proj_state_dict.update(
{
"ff.net.0.proj.weight": state_dict["image_proj"]["proj.0.weight"],
"ff.net.0.proj.bias": state_dict["image_proj"]["proj.0.bias"],
"ff.net.2.weight": state_dict["image_proj"]["proj.2.weight"],
"ff.net.2.bias": state_dict["image_proj"]["proj.2.bias"],
"norm.weight": state_dict["image_proj"]["proj.3.weight"],
"norm.bias": state_dict["image_proj"]["proj.3.bias"],
}
)
image_projection.load_state_dict(image_proj_state_dict)
del image_proj_state_dict
else:
# IP-Adapter Plus
-2
View File
@@ -32,7 +32,6 @@ if is_torch_available():
_import_structure["autoencoder_tiny"] = ["AutoencoderTiny"]
_import_structure["consistency_decoder_vae"] = ["ConsistencyDecoderVAE"]
_import_structure["controlnet"] = ["ControlNetModel"]
_import_structure["controlnetxs"] = ["ControlNetXSModel"]
_import_structure["dual_transformer_2d"] = ["DualTransformer2DModel"]
_import_structure["embeddings"] = ["ImageProjection"]
_import_structure["modeling_utils"] = ["ModelMixin"]
@@ -64,7 +63,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
from .autoencoder_tiny import AutoencoderTiny
from .consistency_decoder_vae import ConsistencyDecoderVAE
from .controlnet import ControlNetModel
from .controlnetxs import ControlNetXSModel
from .dual_transformer_2d import DualTransformer2DModel
from .embeddings import ImageProjection
from .modeling_utils import ModelMixin
+3 -127
View File
@@ -113,14 +113,12 @@ class Attention(nn.Module):
):
super().__init__()
self.inner_dim = out_dim if out_dim is not None else dim_head * heads
self.query_dim = query_dim
self.cross_attention_dim = cross_attention_dim if cross_attention_dim is not None else query_dim
self.upcast_attention = upcast_attention
self.upcast_softmax = upcast_softmax
self.rescale_output_factor = rescale_output_factor
self.residual_connection = residual_connection
self.dropout = dropout
self.fused_projections = False
self.out_dim = out_dim if out_dim is not None else query_dim
# we make use of this private variable to know whether this class is loaded
@@ -182,7 +180,6 @@ class Attention(nn.Module):
else:
linear_cls = LoRACompatibleLinear
self.linear_cls = linear_cls
self.to_q = linear_cls(query_dim, self.inner_dim, bias=bias)
if not self.only_cross_attention:
@@ -695,32 +692,6 @@ class Attention(nn.Module):
return encoder_hidden_states
@torch.no_grad()
def fuse_projections(self, fuse=True):
is_cross_attention = self.cross_attention_dim != self.query_dim
device = self.to_q.weight.data.device
dtype = self.to_q.weight.data.dtype
if not is_cross_attention:
# fetch weight matrices.
concatenated_weights = torch.cat([self.to_q.weight.data, self.to_k.weight.data, self.to_v.weight.data])
in_features = concatenated_weights.shape[1]
out_features = concatenated_weights.shape[0]
# create a new single projection layer and copy over the weights.
self.to_qkv = self.linear_cls(in_features, out_features, bias=False, device=device, dtype=dtype)
self.to_qkv.weight.copy_(concatenated_weights)
else:
concatenated_weights = torch.cat([self.to_k.weight.data, self.to_v.weight.data])
in_features = concatenated_weights.shape[1]
out_features = concatenated_weights.shape[0]
self.to_kv = self.linear_cls(in_features, out_features, bias=False, device=device, dtype=dtype)
self.to_kv.weight.copy_(concatenated_weights)
self.fused_projections = fuse
class AttnProcessor:
r"""
@@ -1213,6 +1184,9 @@ class AttnProcessor2_0:
scale: float = 1.0,
) -> torch.FloatTensor:
residual = hidden_states
args = () if USE_PEFT_BACKEND else (scale,)
if attn.spatial_norm is not None:
hidden_states = attn.spatial_norm(hidden_states, temb)
@@ -1279,103 +1253,6 @@ class AttnProcessor2_0:
return hidden_states
class FusedAttnProcessor2_0:
r"""
Processor for implementing scaled dot-product attention (enabled by default if you're using PyTorch 2.0).
It uses fused projection layers. For self-attention modules, all projection matrices (i.e., query,
key, value) are fused. For cross-attention modules, key and value projection matrices are fused.
<Tip warning={true}>
This API is currently 🧪 experimental in nature and can change in future.
</Tip>
"""
def __init__(self):
if not hasattr(F, "scaled_dot_product_attention"):
raise ImportError(
"FusedAttnProcessor2_0 requires at least PyTorch 2.0, to use it. Please upgrade PyTorch to > 2.0."
)
def __call__(
self,
attn: Attention,
hidden_states: torch.FloatTensor,
encoder_hidden_states: Optional[torch.FloatTensor] = None,
attention_mask: Optional[torch.FloatTensor] = None,
temb: Optional[torch.FloatTensor] = None,
scale: float = 1.0,
) -> torch.FloatTensor:
residual = hidden_states
if attn.spatial_norm is not None:
hidden_states = attn.spatial_norm(hidden_states, temb)
input_ndim = hidden_states.ndim
if input_ndim == 4:
batch_size, channel, height, width = hidden_states.shape
hidden_states = hidden_states.view(batch_size, channel, height * width).transpose(1, 2)
batch_size, sequence_length, _ = (
hidden_states.shape if encoder_hidden_states is None else encoder_hidden_states.shape
)
if attention_mask is not None:
attention_mask = attn.prepare_attention_mask(attention_mask, sequence_length, batch_size)
# scaled_dot_product_attention expects attention_mask shape to be
# (batch, heads, source_length, target_length)
attention_mask = attention_mask.view(batch_size, attn.heads, -1, attention_mask.shape[-1])
if attn.group_norm is not None:
hidden_states = attn.group_norm(hidden_states.transpose(1, 2)).transpose(1, 2)
args = () if USE_PEFT_BACKEND else (scale,)
if encoder_hidden_states is None:
qkv = attn.to_qkv(hidden_states, *args)
split_size = qkv.shape[-1] // 3
query, key, value = torch.split(qkv, split_size, dim=-1)
else:
if attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
query = attn.to_q(hidden_states, *args)
kv = attn.to_kv(encoder_hidden_states, *args)
split_size = kv.shape[-1] // 2
key, value = torch.split(kv, split_size, dim=-1)
inner_dim = key.shape[-1]
head_dim = inner_dim // attn.heads
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
key = key.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
value = value.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
# the output of sdp = (batch, num_heads, seq_len, head_dim)
# TODO: add support for attn.scale when we move to Torch 2.1
hidden_states = F.scaled_dot_product_attention(
query, key, value, attn_mask=attention_mask, dropout_p=0.0, is_causal=False
)
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
hidden_states = hidden_states.to(query.dtype)
# linear proj
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
if input_ndim == 4:
hidden_states = hidden_states.transpose(-1, -2).reshape(batch_size, channel, height, width)
if attn.residual_connection:
hidden_states = hidden_states + residual
hidden_states = hidden_states / attn.rescale_output_factor
return hidden_states
class CustomDiffusionXFormersAttnProcessor(nn.Module):
r"""
Processor for implementing memory efficient attention using xFormers for the Custom Diffusion method.
@@ -2374,7 +2251,6 @@ CROSS_ATTENTION_PROCESSORS = (
AttentionProcessor = Union[
AttnProcessor,
AttnProcessor2_0,
FusedAttnProcessor2_0,
XFormersAttnProcessor,
SlicedAttnProcessor,
AttnAddedKVProcessor,
-39
View File
@@ -22,7 +22,6 @@ from ..utils.accelerate_utils import apply_forward_hook
from .attention_processor import (
ADDED_KV_ATTENTION_PROCESSORS,
CROSS_ATTENTION_PROCESSORS,
Attention,
AttentionProcessor,
AttnAddedKVProcessor,
AttnProcessor,
@@ -449,41 +448,3 @@ class AutoencoderKL(ModelMixin, ConfigMixin, FromOriginalVAEMixin):
return (dec,)
return DecoderOutput(sample=dec)
# Copied from diffusers.models.unet_2d_condition.UNet2DConditionModel.fuse_qkv_projections
def fuse_qkv_projections(self):
"""
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query,
key, value) are fused. For cross-attention modules, key and value projection matrices are fused.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
self.original_attn_processors = None
for _, attn_processor in self.attn_processors.items():
if "Added" in str(attn_processor.__class__.__name__):
raise ValueError("`fuse_qkv_projections()` is not supported for models having added KV projections.")
self.original_attn_processors = self.attn_processors
for module in self.modules():
if isinstance(module, Attention):
module.fuse_projections(fuse=True)
# Copied from diffusers.models.unet_2d_condition.UNet2DConditionModel.unfuse_qkv_projections
def unfuse_qkv_projections(self):
"""Disables the fused QKV projection if enabled.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
if self.original_attn_processors is not None:
self.set_attn_processor(self.original_attn_processors)
-977
View File
@@ -1,977 +0,0 @@
# Copyright 2023 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import math
from dataclasses import dataclass
from typing import Any, Dict, List, Optional, Tuple, Union
import torch
import torch.utils.checkpoint
from torch import nn
from torch.nn import functional as F
from torch.nn.modules.normalization import GroupNorm
from ..configuration_utils import ConfigMixin, register_to_config
from ..utils import BaseOutput, logging
from .attention_processor import (
AttentionProcessor,
)
from .autoencoder_kl import AutoencoderKL
from .lora import LoRACompatibleConv
from .modeling_utils import ModelMixin
from .unet_2d_blocks import (
CrossAttnDownBlock2D,
CrossAttnUpBlock2D,
DownBlock2D,
Downsample2D,
ResnetBlock2D,
Transformer2DModel,
UpBlock2D,
Upsample2D,
)
from .unet_2d_condition import UNet2DConditionModel
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
@dataclass
class ControlNetXSOutput(BaseOutput):
"""
The output of [`ControlNetXSModel`].
Args:
sample (`torch.FloatTensor` of shape `(batch_size, num_channels, height, width)`):
The output of the `ControlNetXSModel`. Unlike `ControlNetOutput` this is NOT to be added to the base model
output, but is already the final output.
"""
sample: torch.FloatTensor = None
# copied from diffusers.models.controlnet.ControlNetConditioningEmbedding
class ControlNetConditioningEmbedding(nn.Module):
"""
Quoting from https://arxiv.org/abs/2302.05543: "Stable Diffusion uses a pre-processing method similar to VQ-GAN
[11] to convert the entire dataset of 512 × 512 images into smaller 64 × 64 latent images for stabilized
training. This requires ControlNets to convert image-based conditions to 64 × 64 feature space to match the
convolution size. We use a tiny network E(·) of four convolution layers with 4 × 4 kernels and 2 × 2 strides
(activated by ReLU, channels are 16, 32, 64, 128, initialized with Gaussian weights, trained jointly with the full
model) to encode image-space conditions ... into feature maps ..."
"""
def __init__(
self,
conditioning_embedding_channels: int,
conditioning_channels: int = 3,
block_out_channels: Tuple[int, ...] = (16, 32, 96, 256),
):
super().__init__()
self.conv_in = nn.Conv2d(conditioning_channels, block_out_channels[0], kernel_size=3, padding=1)
self.blocks = nn.ModuleList([])
for i in range(len(block_out_channels) - 1):
channel_in = block_out_channels[i]
channel_out = block_out_channels[i + 1]
self.blocks.append(nn.Conv2d(channel_in, channel_in, kernel_size=3, padding=1))
self.blocks.append(nn.Conv2d(channel_in, channel_out, kernel_size=3, padding=1, stride=2))
self.conv_out = zero_module(
nn.Conv2d(block_out_channels[-1], conditioning_embedding_channels, kernel_size=3, padding=1)
)
def forward(self, conditioning):
embedding = self.conv_in(conditioning)
embedding = F.silu(embedding)
for block in self.blocks:
embedding = block(embedding)
embedding = F.silu(embedding)
embedding = self.conv_out(embedding)
return embedding
class ControlNetXSModel(ModelMixin, ConfigMixin):
r"""
A ControlNet-XS model
This model inherits from [`ModelMixin`] and [`ConfigMixin`]. Check the superclass documentation for it's generic
methods implemented for all models (such as downloading or saving).
Most of parameters for this model are passed into the [`UNet2DConditionModel`] it creates. Check the documentation
of [`UNet2DConditionModel`] for them.
Parameters:
conditioning_channels (`int`, defaults to 3):
Number of channels of conditioning input (e.g. an image)
controlnet_conditioning_channel_order (`str`, defaults to `"rgb"`):
The channel order of conditional image. Will convert to `rgb` if it's `bgr`.
conditioning_embedding_out_channels (`tuple[int]`, defaults to `(16, 32, 96, 256)`):
The tuple of output channel for each block in the `controlnet_cond_embedding` layer.
time_embedding_input_dim (`int`, defaults to 320):
Dimension of input into time embedding. Needs to be same as in the base model.
time_embedding_dim (`int`, defaults to 1280):
Dimension of output from time embedding. Needs to be same as in the base model.
learn_embedding (`bool`, defaults to `False`):
Whether to use time embedding of the control model. If yes, the time embedding is a linear interpolation of
the time embeddings of the control and base model with interpolation parameter `time_embedding_mix**3`.
time_embedding_mix (`float`, defaults to 1.0):
Linear interpolation parameter used if `learn_embedding` is `True`. A value of 1.0 means only the
control model's time embedding will be used. A value of 0.0 means only the base model's time embedding will be used.
base_model_channel_sizes (`Dict[str, List[Tuple[int]]]`):
Channel sizes of each subblock of base model. Use `gather_subblock_sizes` on your base model to compute it.
"""
@classmethod
def init_original(cls, base_model: UNet2DConditionModel, is_sdxl=True):
"""
Create a ControlNetXS model with the same parameters as in the original paper (https://github.com/vislearn/ControlNet-XS).
Parameters:
base_model (`UNet2DConditionModel`):
Base UNet model. Needs to be either StableDiffusion or StableDiffusion-XL.
is_sdxl (`bool`, defaults to `True`):
Whether passed `base_model` is a StableDiffusion-XL model.
"""
def get_dim_attn_heads(base_model: UNet2DConditionModel, size_ratio: float, num_attn_heads: int):
"""
Currently, diffusers can only set the dimension of attention heads (see https://github.com/huggingface/diffusers/issues/2011#issuecomment-1547958131 for why).
The original ControlNet-XS model, however, define the number of attention heads.
That's why compute the dimensions needed to get the correct number of attention heads.
"""
block_out_channels = [int(size_ratio * c) for c in base_model.config.block_out_channels]
dim_attn_heads = [math.ceil(c / num_attn_heads) for c in block_out_channels]
return dim_attn_heads
if is_sdxl:
return ControlNetXSModel.from_unet(
base_model,
time_embedding_mix=0.95,
learn_embedding=True,
size_ratio=0.1,
conditioning_embedding_out_channels=(16, 32, 96, 256),
num_attention_heads=get_dim_attn_heads(base_model, 0.1, 64),
)
else:
return ControlNetXSModel.from_unet(
base_model,
time_embedding_mix=1.0,
learn_embedding=True,
size_ratio=0.0125,
conditioning_embedding_out_channels=(16, 32, 96, 256),
num_attention_heads=get_dim_attn_heads(base_model, 0.0125, 8),
)
@classmethod
def _gather_subblock_sizes(cls, unet: UNet2DConditionModel, base_or_control: str):
"""To create correctly sized connections between base and control model, we need to know
the input and output channels of each subblock.
Parameters:
unet (`UNet2DConditionModel`):
Unet of which the subblock channels sizes are to be gathered.
base_or_control (`str`):
Needs to be either "base" or "control". If "base", decoder is also considered.
"""
if base_or_control not in ["base", "control"]:
raise ValueError("`base_or_control` needs to be either `base` or `control`")
channel_sizes = {"down": [], "mid": [], "up": []}
# input convolution
channel_sizes["down"].append((unet.conv_in.in_channels, unet.conv_in.out_channels))
# encoder blocks
for module in unet.down_blocks:
if isinstance(module, (CrossAttnDownBlock2D, DownBlock2D)):
for r in module.resnets:
channel_sizes["down"].append((r.in_channels, r.out_channels))
if module.downsamplers:
channel_sizes["down"].append(
(module.downsamplers[0].channels, module.downsamplers[0].out_channels)
)
else:
raise ValueError(f"Encountered unknown module of type {type(module)} while creating ControlNet-XS.")
# middle block
channel_sizes["mid"].append((unet.mid_block.resnets[0].in_channels, unet.mid_block.resnets[0].out_channels))
# decoder blocks
if base_or_control == "base":
for module in unet.up_blocks:
if isinstance(module, (CrossAttnUpBlock2D, UpBlock2D)):
for r in module.resnets:
channel_sizes["up"].append((r.in_channels, r.out_channels))
else:
raise ValueError(
f"Encountered unknown module of type {type(module)} while creating ControlNet-XS."
)
return channel_sizes
@register_to_config
def __init__(
self,
conditioning_channels: int = 3,
conditioning_embedding_out_channels: Tuple[int] = (16, 32, 96, 256),
controlnet_conditioning_channel_order: str = "rgb",
time_embedding_input_dim: int = 320,
time_embedding_dim: int = 1280,
time_embedding_mix: float = 1.0,
learn_embedding: bool = False,
base_model_channel_sizes: Dict[str, List[Tuple[int]]] = {
"down": [
(4, 320),
(320, 320),
(320, 320),
(320, 320),
(320, 640),
(640, 640),
(640, 640),
(640, 1280),
(1280, 1280),
],
"mid": [(1280, 1280)],
"up": [
(2560, 1280),
(2560, 1280),
(1920, 1280),
(1920, 640),
(1280, 640),
(960, 640),
(960, 320),
(640, 320),
(640, 320),
],
},
sample_size: Optional[int] = None,
down_block_types: Tuple[str] = (
"CrossAttnDownBlock2D",
"CrossAttnDownBlock2D",
"CrossAttnDownBlock2D",
"DownBlock2D",
),
up_block_types: Tuple[str] = ("UpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D"),
block_out_channels: Tuple[int] = (320, 640, 1280, 1280),
norm_num_groups: Optional[int] = 32,
cross_attention_dim: Union[int, Tuple[int]] = 1280,
transformer_layers_per_block: Union[int, Tuple[int], Tuple[Tuple]] = 1,
num_attention_heads: Optional[Union[int, Tuple[int]]] = 8,
upcast_attention: bool = False,
):
super().__init__()
# 1 - Create control unet
self.control_model = UNet2DConditionModel(
sample_size=sample_size,
down_block_types=down_block_types,
up_block_types=up_block_types,
block_out_channels=block_out_channels,
norm_num_groups=norm_num_groups,
cross_attention_dim=cross_attention_dim,
transformer_layers_per_block=transformer_layers_per_block,
attention_head_dim=num_attention_heads,
use_linear_projection=True,
upcast_attention=upcast_attention,
time_embedding_dim=time_embedding_dim,
)
# 2 - Do model surgery on control model
# 2.1 - Allow to use the same time information as the base model
adjust_time_dims(self.control_model, time_embedding_input_dim, time_embedding_dim)
# 2.2 - Allow for information infusion from base model
# We concat the output of each base encoder subblocks to the input of the next control encoder subblock
# (We ignore the 1st element, as it represents the `conv_in`.)
extra_input_channels = [input_channels for input_channels, _ in base_model_channel_sizes["down"][1:]]
it_extra_input_channels = iter(extra_input_channels)
for b, block in enumerate(self.control_model.down_blocks):
for r in range(len(block.resnets)):
increase_block_input_in_encoder_resnet(
self.control_model, block_no=b, resnet_idx=r, by=next(it_extra_input_channels)
)
if block.downsamplers:
increase_block_input_in_encoder_downsampler(
self.control_model, block_no=b, by=next(it_extra_input_channels)
)
increase_block_input_in_mid_resnet(self.control_model, by=extra_input_channels[-1])
# 2.3 - Make group norms work with modified channel sizes
adjust_group_norms(self.control_model)
# 3 - Gather Channel Sizes
self.ch_inout_ctrl = ControlNetXSModel._gather_subblock_sizes(self.control_model, base_or_control="control")
self.ch_inout_base = base_model_channel_sizes
# 4 - Build connections between base and control model
self.down_zero_convs_out = nn.ModuleList([])
self.down_zero_convs_in = nn.ModuleList([])
self.middle_block_out = nn.ModuleList([])
self.middle_block_in = nn.ModuleList([])
self.up_zero_convs_out = nn.ModuleList([])
self.up_zero_convs_in = nn.ModuleList([])
for ch_io_base in self.ch_inout_base["down"]:
self.down_zero_convs_in.append(self._make_zero_conv(in_channels=ch_io_base[1], out_channels=ch_io_base[1]))
for i in range(len(self.ch_inout_ctrl["down"])):
self.down_zero_convs_out.append(
self._make_zero_conv(self.ch_inout_ctrl["down"][i][1], self.ch_inout_base["down"][i][1])
)
self.middle_block_out = self._make_zero_conv(
self.ch_inout_ctrl["mid"][-1][1], self.ch_inout_base["mid"][-1][1]
)
self.up_zero_convs_out.append(
self._make_zero_conv(self.ch_inout_ctrl["down"][-1][1], self.ch_inout_base["mid"][-1][1])
)
for i in range(1, len(self.ch_inout_ctrl["down"])):
self.up_zero_convs_out.append(
self._make_zero_conv(self.ch_inout_ctrl["down"][-(i + 1)][1], self.ch_inout_base["up"][i - 1][1])
)
# 5 - Create conditioning hint embedding
self.controlnet_cond_embedding = ControlNetConditioningEmbedding(
conditioning_embedding_channels=block_out_channels[0],
block_out_channels=conditioning_embedding_out_channels,
conditioning_channels=conditioning_channels,
)
# In the mininal implementation setting, we only need the control model up to the mid block
del self.control_model.up_blocks
del self.control_model.conv_norm_out
del self.control_model.conv_out
@classmethod
def from_unet(
cls,
unet: UNet2DConditionModel,
conditioning_channels: int = 3,
conditioning_embedding_out_channels: Tuple[int] = (16, 32, 96, 256),
controlnet_conditioning_channel_order: str = "rgb",
learn_embedding: bool = False,
time_embedding_mix: float = 1.0,
block_out_channels: Optional[Tuple[int]] = None,
size_ratio: Optional[float] = None,
num_attention_heads: Optional[Union[int, Tuple[int]]] = 8,
norm_num_groups: Optional[int] = None,
):
r"""
Instantiate a [`ControlNetXSModel`] from [`UNet2DConditionModel`].
Parameters:
unet (`UNet2DConditionModel`):
The UNet model we want to control. The dimensions of the ControlNetXSModel will be adapted to it.
conditioning_channels (`int`, defaults to 3):
Number of channels of conditioning input (e.g. an image)
conditioning_embedding_out_channels (`tuple[int]`, defaults to `(16, 32, 96, 256)`):
The tuple of output channel for each block in the `controlnet_cond_embedding` layer.
controlnet_conditioning_channel_order (`str`, defaults to `"rgb"`):
The channel order of conditional image. Will convert to `rgb` if it's `bgr`.
learn_embedding (`bool`, defaults to `False`):
Wether to use time embedding of the control model. If yes, the time embedding is a linear interpolation
of the time embeddings of the control and base model with interpolation parameter
`time_embedding_mix**3`.
time_embedding_mix (`float`, defaults to 1.0):
Linear interpolation parameter used if `learn_embedding` is `True`.
block_out_channels (`Tuple[int]`, *optional*):
Down blocks output channels in control model. Either this or `size_ratio` must be given.
size_ratio (float, *optional*):
When given, block_out_channels is set to a relative fraction of the base model's block_out_channels.
Either this or `block_out_channels` must be given.
num_attention_heads (`Union[int, Tuple[int]]`, *optional*):
The dimension of the attention heads. The naming seems a bit confusing and it is, see https://github.com/huggingface/diffusers/issues/2011#issuecomment-1547958131 for why.
norm_num_groups (int, *optional*, defaults to `None`):
The number of groups to use for the normalization of the control unet. If `None`,
`int(unet.config.norm_num_groups * size_ratio)` is taken.
"""
# Check input
fixed_size = block_out_channels is not None
relative_size = size_ratio is not None
if not (fixed_size ^ relative_size):
raise ValueError(
"Pass exactly one of `block_out_channels` (for absolute sizing) or `control_model_ratio` (for relative sizing)."
)
# Create model
if block_out_channels is None:
block_out_channels = [int(size_ratio * c) for c in unet.config.block_out_channels]
# Check that attention heads and group norms match channel sizes
# - attention heads
def attn_heads_match_channel_sizes(attn_heads, channel_sizes):
if isinstance(attn_heads, (tuple, list)):
return all(c % a == 0 for a, c in zip(attn_heads, channel_sizes))
else:
return all(c % attn_heads == 0 for c in channel_sizes)
num_attention_heads = num_attention_heads or unet.config.attention_head_dim
if not attn_heads_match_channel_sizes(num_attention_heads, block_out_channels):
raise ValueError(
f"The dimension of attention heads ({num_attention_heads}) must divide `block_out_channels` ({block_out_channels}). If you didn't set `num_attention_heads` the default settings don't match your model. Set `num_attention_heads` manually."
)
# - group norms
def group_norms_match_channel_sizes(num_groups, channel_sizes):
return all(c % num_groups == 0 for c in channel_sizes)
if norm_num_groups is None:
if group_norms_match_channel_sizes(unet.config.norm_num_groups, block_out_channels):
norm_num_groups = unet.config.norm_num_groups
else:
norm_num_groups = min(block_out_channels)
if group_norms_match_channel_sizes(norm_num_groups, block_out_channels):
print(
f"`norm_num_groups` was set to `min(block_out_channels)` (={norm_num_groups}) so it divides all block_out_channels` ({block_out_channels}). Set it explicitly to remove this information."
)
else:
raise ValueError(
f"`block_out_channels` ({block_out_channels}) don't match the base models `norm_num_groups` ({unet.config.norm_num_groups}). Setting `norm_num_groups` to `min(block_out_channels)` ({norm_num_groups}) didn't fix this. Pass `norm_num_groups` explicitly so it divides all block_out_channels."
)
def get_time_emb_input_dim(unet: UNet2DConditionModel):
return unet.time_embedding.linear_1.in_features
def get_time_emb_dim(unet: UNet2DConditionModel):
return unet.time_embedding.linear_2.out_features
# Clone params from base unet if
# (i) it's required to build SD or SDXL, and
# (ii) it's not used for the time embedding (as time embedding of control model is never used), and
# (iii) it's not set further below anyway
to_keep = [
"cross_attention_dim",
"down_block_types",
"sample_size",
"transformer_layers_per_block",
"up_block_types",
"upcast_attention",
]
kwargs = {k: v for k, v in dict(unet.config).items() if k in to_keep}
kwargs.update(block_out_channels=block_out_channels)
kwargs.update(num_attention_heads=num_attention_heads)
kwargs.update(norm_num_groups=norm_num_groups)
# Add controlnetxs-specific params
kwargs.update(
conditioning_channels=conditioning_channels,
controlnet_conditioning_channel_order=controlnet_conditioning_channel_order,
time_embedding_input_dim=get_time_emb_input_dim(unet),
time_embedding_dim=get_time_emb_dim(unet),
time_embedding_mix=time_embedding_mix,
learn_embedding=learn_embedding,
base_model_channel_sizes=ControlNetXSModel._gather_subblock_sizes(unet, base_or_control="base"),
conditioning_embedding_out_channels=conditioning_embedding_out_channels,
)
return cls(**kwargs)
@property
def attn_processors(self) -> Dict[str, AttentionProcessor]:
r"""
Returns:
`dict` of attention processors: A dictionary containing all attention processors used in the model with
indexed by its weight name.
"""
return self.control_model.attn_processors
def set_attn_processor(
self, processor: Union[AttentionProcessor, Dict[str, AttentionProcessor]], _remove_lora=False
):
r"""
Sets the attention processor to use to compute attention.
Parameters:
processor (`dict` of `AttentionProcessor` or only `AttentionProcessor`):
The instantiated processor class or a dictionary of processor classes that will be set as the processor
for **all** `Attention` layers.
If `processor` is a dict, the key needs to define the path to the corresponding cross attention
processor. This is strongly recommended when setting trainable attention processors.
"""
self.control_model.set_attn_processor(processor, _remove_lora)
def set_default_attn_processor(self):
"""
Disables custom attention processors and sets the default attention implementation.
"""
self.control_model.set_default_attn_processor()
def set_attention_slice(self, slice_size):
r"""
Enable sliced attention computation.
When this option is enabled, the attention module splits the input tensor in slices to compute attention in
several steps. This is useful for saving some memory in exchange for a small decrease in speed.
Args:
slice_size (`str` or `int` or `list(int)`, *optional*, defaults to `"auto"`):
When `"auto"`, input to the attention heads is halved, so attention is computed in two steps. If
`"max"`, maximum amount of memory is saved by running only one slice at a time. If a number is
provided, uses as many slices as `attention_head_dim // slice_size`. In this case, `attention_head_dim`
must be a multiple of `slice_size`.
"""
self.control_model.set_attention_slice(slice_size)
def _set_gradient_checkpointing(self, module, value=False):
if isinstance(module, (UNet2DConditionModel)):
if value:
module.enable_gradient_checkpointing()
else:
module.disable_gradient_checkpointing()
def forward(
self,
base_model: UNet2DConditionModel,
sample: torch.FloatTensor,
timestep: Union[torch.Tensor, float, int],
encoder_hidden_states: torch.Tensor,
controlnet_cond: torch.Tensor,
conditioning_scale: float = 1.0,
class_labels: Optional[torch.Tensor] = None,
timestep_cond: Optional[torch.Tensor] = None,
attention_mask: Optional[torch.Tensor] = None,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
added_cond_kwargs: Optional[Dict[str, torch.Tensor]] = None,
return_dict: bool = True,
) -> Union[ControlNetXSOutput, Tuple]:
"""
The [`ControlNetModel`] forward method.
Args:
base_model (`UNet2DConditionModel`):
The base unet model we want to control.
sample (`torch.FloatTensor`):
The noisy input tensor.
timestep (`Union[torch.Tensor, float, int]`):
The number of timesteps to denoise an input.
encoder_hidden_states (`torch.Tensor`):
The encoder hidden states.
controlnet_cond (`torch.FloatTensor`):
The conditional input tensor of shape `(batch_size, sequence_length, hidden_size)`.
conditioning_scale (`float`, defaults to `1.0`):
How much the control model affects the base model outputs.
class_labels (`torch.Tensor`, *optional*, defaults to `None`):
Optional class labels for conditioning. Their embeddings will be summed with the timestep embeddings.
timestep_cond (`torch.Tensor`, *optional*, defaults to `None`):
Additional conditional embeddings for timestep. If provided, the embeddings will be summed with the
timestep_embedding passed through the `self.time_embedding` layer to obtain the final timestep
embeddings.
attention_mask (`torch.Tensor`, *optional*, defaults to `None`):
An attention mask of shape `(batch, key_tokens)` is applied to `encoder_hidden_states`. If `1` the mask
is kept, otherwise if `0` it is discarded. Mask will be converted into a bias, which adds large
negative values to the attention scores corresponding to "discard" tokens.
added_cond_kwargs (`dict`):
Additional conditions for the Stable Diffusion XL UNet.
cross_attention_kwargs (`dict[str]`, *optional*, defaults to `None`):
A kwargs dictionary that if specified is passed along to the `AttnProcessor`.
return_dict (`bool`, defaults to `True`):
Whether or not to return a [`~models.controlnet.ControlNetOutput`] instead of a plain tuple.
Returns:
[`~models.controlnetxs.ControlNetXSOutput`] **or** `tuple`:
If `return_dict` is `True`, a [`~models.controlnetxs.ControlNetXSOutput`] is returned, otherwise a
tuple is returned where the first element is the sample tensor.
"""
# check channel order
channel_order = self.config.controlnet_conditioning_channel_order
if channel_order == "rgb":
# in rgb order by default
...
elif channel_order == "bgr":
controlnet_cond = torch.flip(controlnet_cond, dims=[1])
else:
raise ValueError(f"unknown `controlnet_conditioning_channel_order`: {channel_order}")
# scale control strength
n_connections = len(self.down_zero_convs_out) + 1 + len(self.up_zero_convs_out)
scale_list = torch.full((n_connections,), conditioning_scale)
# prepare attention_mask
if attention_mask is not None:
attention_mask = (1 - attention_mask.to(sample.dtype)) * -10000.0
attention_mask = attention_mask.unsqueeze(1)
# 1. time
timesteps = timestep
if not torch.is_tensor(timesteps):
# TODO: this requires sync between CPU and GPU. So try to pass timesteps as tensors if you can
# This would be a good case for the `match` statement (Python 3.10+)
is_mps = sample.device.type == "mps"
if isinstance(timestep, float):
dtype = torch.float32 if is_mps else torch.float64
else:
dtype = torch.int32 if is_mps else torch.int64
timesteps = torch.tensor([timesteps], dtype=dtype, device=sample.device)
elif len(timesteps.shape) == 0:
timesteps = timesteps[None].to(sample.device)
# broadcast to batch dimension in a way that's compatible with ONNX/Core ML
timesteps = timesteps.expand(sample.shape[0])
t_emb = base_model.time_proj(timesteps)
# timesteps does not contain any weights and will always return f32 tensors
# but time_embedding might actually be running in fp16. so we need to cast here.
# there might be better ways to encapsulate this.
t_emb = t_emb.to(dtype=sample.dtype)
if self.config.learn_embedding:
ctrl_temb = self.control_model.time_embedding(t_emb, timestep_cond)
base_temb = base_model.time_embedding(t_emb, timestep_cond)
interpolation_param = self.config.time_embedding_mix**0.3
temb = ctrl_temb * interpolation_param + base_temb * (1 - interpolation_param)
else:
temb = base_model.time_embedding(t_emb)
# added time & text embeddings
aug_emb = None
if base_model.class_embedding is not None:
if class_labels is None:
raise ValueError("class_labels should be provided when num_class_embeds > 0")
if base_model.config.class_embed_type == "timestep":
class_labels = base_model.time_proj(class_labels)
class_emb = base_model.class_embedding(class_labels).to(dtype=self.dtype)
temb = temb + class_emb
if base_model.config.addition_embed_type is not None:
if base_model.config.addition_embed_type == "text":
aug_emb = base_model.add_embedding(encoder_hidden_states)
elif base_model.config.addition_embed_type == "text_image":
raise NotImplementedError()
elif base_model.config.addition_embed_type == "text_time":
# SDXL - style
if "text_embeds" not in added_cond_kwargs:
raise ValueError(
f"{self.__class__} has the config param `addition_embed_type` set to 'text_time' which requires the keyword argument `text_embeds` to be passed in `added_cond_kwargs`"
)
text_embeds = added_cond_kwargs.get("text_embeds")
if "time_ids" not in added_cond_kwargs:
raise ValueError(
f"{self.__class__} has the config param `addition_embed_type` set to 'text_time' which requires the keyword argument `time_ids` to be passed in `added_cond_kwargs`"
)
time_ids = added_cond_kwargs.get("time_ids")
time_embeds = base_model.add_time_proj(time_ids.flatten())
time_embeds = time_embeds.reshape((text_embeds.shape[0], -1))
add_embeds = torch.concat([text_embeds, time_embeds], dim=-1)
add_embeds = add_embeds.to(temb.dtype)
aug_emb = base_model.add_embedding(add_embeds)
elif base_model.config.addition_embed_type == "image":
raise NotImplementedError()
elif base_model.config.addition_embed_type == "image_hint":
raise NotImplementedError()
temb = temb + aug_emb if aug_emb is not None else temb
# text embeddings
cemb = encoder_hidden_states
# Preparation
guided_hint = self.controlnet_cond_embedding(controlnet_cond)
h_ctrl = h_base = sample
hs_base, hs_ctrl = [], []
it_down_convs_in, it_down_convs_out, it_dec_convs_in, it_up_convs_out = map(
iter, (self.down_zero_convs_in, self.down_zero_convs_out, self.up_zero_convs_in, self.up_zero_convs_out)
)
scales = iter(scale_list)
base_down_subblocks = to_sub_blocks(base_model.down_blocks)
ctrl_down_subblocks = to_sub_blocks(self.control_model.down_blocks)
base_mid_subblocks = to_sub_blocks([base_model.mid_block])
ctrl_mid_subblocks = to_sub_blocks([self.control_model.mid_block])
base_up_subblocks = to_sub_blocks(base_model.up_blocks)
# Cross Control
# 0 - conv in
h_base = base_model.conv_in(h_base)
h_ctrl = self.control_model.conv_in(h_ctrl)
if guided_hint is not None:
h_ctrl += guided_hint
h_base = h_base + next(it_down_convs_out)(h_ctrl) * next(scales) # D - add ctrl -> base
hs_base.append(h_base)
hs_ctrl.append(h_ctrl)
# 1 - down
for m_base, m_ctrl in zip(base_down_subblocks, ctrl_down_subblocks):
h_ctrl = torch.cat([h_ctrl, next(it_down_convs_in)(h_base)], dim=1) # A - concat base -> ctrl
h_base = m_base(h_base, temb, cemb, attention_mask, cross_attention_kwargs) # B - apply base subblock
h_ctrl = m_ctrl(h_ctrl, temb, cemb, attention_mask, cross_attention_kwargs) # C - apply ctrl subblock
h_base = h_base + next(it_down_convs_out)(h_ctrl) * next(scales) # D - add ctrl -> base
hs_base.append(h_base)
hs_ctrl.append(h_ctrl)
# 2 - mid
h_ctrl = torch.cat([h_ctrl, next(it_down_convs_in)(h_base)], dim=1) # A - concat base -> ctrl
for m_base, m_ctrl in zip(base_mid_subblocks, ctrl_mid_subblocks):
h_base = m_base(h_base, temb, cemb, attention_mask, cross_attention_kwargs) # B - apply base subblock
h_ctrl = m_ctrl(h_ctrl, temb, cemb, attention_mask, cross_attention_kwargs) # C - apply ctrl subblock
h_base = h_base + self.middle_block_out(h_ctrl) * next(scales) # D - add ctrl -> base
# 3 - up
for i, m_base in enumerate(base_up_subblocks):
h_base = h_base + next(it_up_convs_out)(hs_ctrl.pop()) * next(scales) # add info from ctrl encoder
h_base = torch.cat([h_base, hs_base.pop()], dim=1) # concat info from base encoder+ctrl encoder
h_base = m_base(h_base, temb, cemb, attention_mask, cross_attention_kwargs)
h_base = base_model.conv_norm_out(h_base)
h_base = base_model.conv_act(h_base)
h_base = base_model.conv_out(h_base)
if not return_dict:
return h_base
return ControlNetXSOutput(sample=h_base)
def _make_zero_conv(self, in_channels, out_channels=None):
# keep running track of channels sizes
self.in_channels = in_channels
self.out_channels = out_channels or in_channels
return zero_module(nn.Conv2d(in_channels, out_channels, 1, padding=0))
@torch.no_grad()
def _check_if_vae_compatible(self, vae: AutoencoderKL):
condition_downscale_factor = 2 ** (len(self.config.conditioning_embedding_out_channels) - 1)
vae_downscale_factor = 2 ** (len(vae.config.block_out_channels) - 1)
compatible = condition_downscale_factor == vae_downscale_factor
return compatible, condition_downscale_factor, vae_downscale_factor
class SubBlock(nn.ModuleList):
"""A SubBlock is the largest piece of either base or control model, that is executed independently of the other model respectively.
Before each subblock, information is concatted from base to control. And after each subblock, information is added from control to base.
"""
def __init__(self, ms, *args, **kwargs):
if not is_iterable(ms):
ms = [ms]
super().__init__(ms, *args, **kwargs)
def forward(
self,
x: torch.Tensor,
temb: torch.Tensor,
cemb: torch.Tensor,
attention_mask: Optional[torch.Tensor] = None,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
):
"""Iterate through children and pass correct information to each."""
for m in self:
if isinstance(m, ResnetBlock2D):
x = m(x, temb)
elif isinstance(m, Transformer2DModel):
x = m(x, cemb, attention_mask=attention_mask, cross_attention_kwargs=cross_attention_kwargs).sample
elif isinstance(m, Downsample2D):
x = m(x)
elif isinstance(m, Upsample2D):
x = m(x)
else:
raise ValueError(
f"Type of m is {type(m)} but should be `ResnetBlock2D`, `Transformer2DModel`, `Downsample2D` or `Upsample2D`"
)
return x
def adjust_time_dims(unet: UNet2DConditionModel, in_dim: int, out_dim: int):
unet.time_embedding.linear_1 = nn.Linear(in_dim, out_dim)
def increase_block_input_in_encoder_resnet(unet: UNet2DConditionModel, block_no, resnet_idx, by):
"""Increase channels sizes to allow for additional concatted information from base model"""
r = unet.down_blocks[block_no].resnets[resnet_idx]
old_norm1, old_conv1 = r.norm1, r.conv1
# norm
norm_args = "num_groups num_channels eps affine".split(" ")
for a in norm_args:
assert hasattr(old_norm1, a)
norm_kwargs = {a: getattr(old_norm1, a) for a in norm_args}
norm_kwargs["num_channels"] += by # surgery done here
# conv1
conv1_args = (
"in_channels out_channels kernel_size stride padding dilation groups bias padding_mode lora_layer".split(" ")
)
for a in conv1_args:
assert hasattr(old_conv1, a)
conv1_kwargs = {a: getattr(old_conv1, a) for a in conv1_args}
conv1_kwargs["bias"] = "bias" in conv1_kwargs # as param, bias is a boolean, but as attr, it's a tensor.
conv1_kwargs["in_channels"] += by # surgery done here
# conv_shortcut
# as we changed the input size of the block, the input and output sizes are likely different,
# therefore we need a conv_shortcut (simply adding won't work)
conv_shortcut_args_kwargs = {
"in_channels": conv1_kwargs["in_channels"],
"out_channels": conv1_kwargs["out_channels"],
# default arguments from resnet.__init__
"kernel_size": 1,
"stride": 1,
"padding": 0,
"bias": True,
}
# swap old with new modules
unet.down_blocks[block_no].resnets[resnet_idx].norm1 = GroupNorm(**norm_kwargs)
unet.down_blocks[block_no].resnets[resnet_idx].conv1 = LoRACompatibleConv(**conv1_kwargs)
unet.down_blocks[block_no].resnets[resnet_idx].conv_shortcut = LoRACompatibleConv(**conv_shortcut_args_kwargs)
unet.down_blocks[block_no].resnets[resnet_idx].in_channels += by # surgery done here
def increase_block_input_in_encoder_downsampler(unet: UNet2DConditionModel, block_no, by):
"""Increase channels sizes to allow for additional concatted information from base model"""
old_down = unet.down_blocks[block_no].downsamplers[0].conv
# conv1
args = "in_channels out_channels kernel_size stride padding dilation groups bias padding_mode lora_layer".split(
" "
)
for a in args:
assert hasattr(old_down, a)
kwargs = {a: getattr(old_down, a) for a in args}
kwargs["bias"] = "bias" in kwargs # as param, bias is a boolean, but as attr, it's a tensor.
kwargs["in_channels"] += by # surgery done here
# swap old with new modules
unet.down_blocks[block_no].downsamplers[0].conv = LoRACompatibleConv(**kwargs)
unet.down_blocks[block_no].downsamplers[0].channels += by # surgery done here
def increase_block_input_in_mid_resnet(unet: UNet2DConditionModel, by):
"""Increase channels sizes to allow for additional concatted information from base model"""
m = unet.mid_block.resnets[0]
old_norm1, old_conv1 = m.norm1, m.conv1
# norm
norm_args = "num_groups num_channels eps affine".split(" ")
for a in norm_args:
assert hasattr(old_norm1, a)
norm_kwargs = {a: getattr(old_norm1, a) for a in norm_args}
norm_kwargs["num_channels"] += by # surgery done here
# conv1
conv1_args = (
"in_channels out_channels kernel_size stride padding dilation groups bias padding_mode lora_layer".split(" ")
)
for a in conv1_args:
assert hasattr(old_conv1, a)
conv1_kwargs = {a: getattr(old_conv1, a) for a in conv1_args}
conv1_kwargs["bias"] = "bias" in conv1_kwargs # as param, bias is a boolean, but as attr, it's a tensor.
conv1_kwargs["in_channels"] += by # surgery done here
# conv_shortcut
# as we changed the input size of the block, the input and output sizes are likely different,
# therefore we need a conv_shortcut (simply adding won't work)
conv_shortcut_args_kwargs = {
"in_channels": conv1_kwargs["in_channels"],
"out_channels": conv1_kwargs["out_channels"],
# default arguments from resnet.__init__
"kernel_size": 1,
"stride": 1,
"padding": 0,
"bias": True,
}
# swap old with new modules
unet.mid_block.resnets[0].norm1 = GroupNorm(**norm_kwargs)
unet.mid_block.resnets[0].conv1 = LoRACompatibleConv(**conv1_kwargs)
unet.mid_block.resnets[0].conv_shortcut = LoRACompatibleConv(**conv_shortcut_args_kwargs)
unet.mid_block.resnets[0].in_channels += by # surgery done here
def adjust_group_norms(unet: UNet2DConditionModel, max_num_group: int = 32):
def find_denominator(number, start):
if start >= number:
return number
while start != 0:
residual = number % start
if residual == 0:
return start
start -= 1
for block in [*unet.down_blocks, unet.mid_block]:
# resnets
for r in block.resnets:
if r.norm1.num_groups < max_num_group:
r.norm1.num_groups = find_denominator(r.norm1.num_channels, start=max_num_group)
if r.norm2.num_groups < max_num_group:
r.norm2.num_groups = find_denominator(r.norm2.num_channels, start=max_num_group)
# transformers
if hasattr(block, "attentions"):
for a in block.attentions:
if a.norm.num_groups < max_num_group:
a.norm.num_groups = find_denominator(a.norm.num_channels, start=max_num_group)
def is_iterable(o):
if isinstance(o, str):
return False
try:
iter(o)
return True
except TypeError:
return False
def to_sub_blocks(blocks):
if not is_iterable(blocks):
blocks = [blocks]
sub_blocks = []
for b in blocks:
if hasattr(b, "resnets"):
if hasattr(b, "attentions") and b.attentions is not None:
for r, a in zip(b.resnets, b.attentions):
sub_blocks.append([r, a])
num_resnets = len(b.resnets)
num_attns = len(b.attentions)
if num_resnets > num_attns:
# we can have more resnets than attentions, so add each resnet as separate subblock
for i in range(num_attns, num_resnets):
sub_blocks.append([b.resnets[i]])
else:
for r in b.resnets:
sub_blocks.append([r])
# upsamplers are part of the same subblock
if hasattr(b, "upsamplers") and b.upsamplers is not None:
for u in b.upsamplers:
sub_blocks[-1].extend([u])
# downsamplers are own subblock
if hasattr(b, "downsamplers") and b.downsamplers is not None:
for d in b.downsamplers:
sub_blocks.append([d])
return list(map(SubBlock, sub_blocks))
def zero_module(module):
for p in module.parameters():
nn.init.zeros_(p)
return module
-12
View File
@@ -461,18 +461,6 @@ class ImageProjection(nn.Module):
return image_embeds
class MLPProjection(nn.Module):
def __init__(self, image_embed_dim=1024, cross_attention_dim=1024):
super().__init__()
from .attention import FeedForward
self.ff = FeedForward(image_embed_dim, cross_attention_dim, mult=1, activation_fn="gelu")
self.norm = nn.LayerNorm(cross_attention_dim)
def forward(self, image_embeds: torch.FloatTensor):
return self.norm(self.ff(image_embeds))
class CombinedTimestepLabelEmbeddings(nn.Module):
def __init__(self, num_classes, embedding_dim, class_dropout_prob=0.1):
super().__init__()
-37
View File
@@ -25,7 +25,6 @@ from .activations import get_activation
from .attention_processor import (
ADDED_KV_ATTENTION_PROCESSORS,
CROSS_ATTENTION_PROCESSORS,
Attention,
AttentionProcessor,
AttnAddedKVProcessor,
AttnProcessor,
@@ -795,42 +794,6 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin, UNet2DConditionLoadersMixin)
if hasattr(upsample_block, k) or getattr(upsample_block, k, None) is not None:
setattr(upsample_block, k, None)
def fuse_qkv_projections(self):
"""
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query,
key, value) are fused. For cross-attention modules, key and value projection matrices are fused.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
self.original_attn_processors = None
for _, attn_processor in self.attn_processors.items():
if "Added" in str(attn_processor.__class__.__name__):
raise ValueError("`fuse_qkv_projections()` is not supported for models having added KV projections.")
self.original_attn_processors = self.attn_processors
for module in self.modules():
if isinstance(module, Attention):
module.fuse_projections(fuse=True)
def unfuse_qkv_projections(self):
"""Disables the fused QKV projection if enabled.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
if self.original_attn_processors is not None:
self.set_attn_processor(self.original_attn_processors)
def forward(
self,
sample: torch.FloatTensor,
-11
View File
@@ -19,7 +19,6 @@ from ..utils import (
_dummy_objects = {}
_import_structure = {
"controlnet": [],
"controlnet_xs": [],
"latent_diffusion": [],
"stable_diffusion": [],
"stable_diffusion_xl": [],
@@ -94,12 +93,6 @@ else:
"StableDiffusionXLControlNetPipeline",
]
)
_import_structure["controlnet_xs"].extend(
[
"StableDiffusionControlNetXSPipeline",
"StableDiffusionXLControlNetXSPipeline",
]
)
_import_structure["deepfloyd_if"] = [
"IFImg2ImgPipeline",
"IFImg2ImgSuperResolutionPipeline",
@@ -354,10 +347,6 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
StableDiffusionXLControlNetInpaintPipeline,
StableDiffusionXLControlNetPipeline,
)
from .controlnet_xs import (
StableDiffusionControlNetXSPipeline,
StableDiffusionXLControlNetXSPipeline,
)
from .deepfloyd_if import (
IFImg2ImgPipeline,
IFImg2ImgSuperResolutionPipeline,
@@ -1,68 +0,0 @@
from typing import TYPE_CHECKING
from ...utils import (
DIFFUSERS_SLOW_IMPORT,
OptionalDependencyNotAvailable,
_LazyModule,
get_objects_from_module,
is_flax_available,
is_torch_available,
is_transformers_available,
)
_dummy_objects = {}
_import_structure = {}
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_torch_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_torch_and_transformers_objects))
else:
_import_structure["pipeline_controlnet_xs"] = ["StableDiffusionControlNetXSPipeline"]
_import_structure["pipeline_controlnet_xs_sd_xl"] = ["StableDiffusionXLControlNetXSPipeline"]
try:
if not (is_transformers_available() and is_flax_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_flax_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_flax_and_transformers_objects))
else:
pass # _import_structure["pipeline_flax_controlnet"] = ["FlaxStableDiffusionControlNetPipeline"]
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_torch_and_transformers_objects import *
else:
from .pipeline_controlnet_xs import StableDiffusionControlNetXSPipeline
from .pipeline_controlnet_xs_sd_xl import StableDiffusionXLControlNetXSPipeline
try:
if not (is_transformers_available() and is_flax_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_flax_and_transformers_objects import * # noqa F403
else:
pass # from .pipeline_flax_controlnet import FlaxStableDiffusionControlNetPipeline
else:
import sys
sys.modules[__name__] = _LazyModule(
__name__,
globals()["__file__"],
_import_structure,
module_spec=__spec__,
)
for name, value in _dummy_objects.items():
setattr(sys.modules[__name__], name, value)
@@ -1,944 +0,0 @@
# Copyright 2023 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import inspect
from typing import Any, Callable, Dict, List, Optional, Union
import numpy as np
import PIL.Image
import torch
import torch.nn.functional as F
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from ...image_processor import PipelineImageInput, VaeImageProcessor
from ...loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from ...models import AutoencoderKL, ControlNetXSModel, UNet2DConditionModel
from ...models.lora import adjust_lora_scale_text_encoder
from ...schedulers import KarrasDiffusionSchedulers
from ...utils import (
USE_PEFT_BACKEND,
deprecate,
logging,
replace_example_docstring,
scale_lora_layers,
unscale_lora_layers,
)
from ...utils.torch_utils import is_compiled_module, is_torch_version, randn_tensor
from ..pipeline_utils import DiffusionPipeline
from ..stable_diffusion.pipeline_output import StableDiffusionPipelineOutput
from ..stable_diffusion.safety_checker import StableDiffusionSafetyChecker
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
EXAMPLE_DOC_STRING = """
Examples:
```py
>>> # !pip install opencv-python transformers accelerate
>>> from diffusers import StableDiffusionControlNetXSPipeline, ControlNetXSModel
>>> from diffusers.utils import load_image
>>> import numpy as np
>>> import torch
>>> import cv2
>>> from PIL import Image
>>> prompt = "aerial view, a futuristic research complex in a bright foggy jungle, hard lighting"
>>> negative_prompt = "low quality, bad quality, sketches"
>>> # download an image
>>> image = load_image(
... "https://hf.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/hf-logo.png"
... )
>>> # initialize the models and pipeline
>>> controlnet_conditioning_scale = 0.5
>>> controlnet = ControlNetXSModel.from_pretrained(
... "UmerHA/ConrolNetXS-SD2.1-canny", torch_dtype=torch.float16
... )
>>> pipe = StableDiffusionControlNetXSPipeline.from_pretrained(
... "stabilityai/stable-diffusion-2-1", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> pipe.enable_model_cpu_offload()
>>> # get canny image
>>> image = np.array(image)
>>> image = cv2.Canny(image, 100, 200)
>>> image = image[:, :, None]
>>> image = np.concatenate([image, image, image], axis=2)
>>> canny_image = Image.fromarray(image)
>>> # generate image
>>> image = pipe(
... prompt, controlnet_conditioning_scale=controlnet_conditioning_scale, image=canny_image
... ).images[0]
```
"""
class StableDiffusionControlNetXSPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image generation using Stable Diffusion with ControlNet-XS guidance.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
text_encoder ([`~transformers.CLIPTextModel`]):
Frozen text-encoder ([clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14)).
tokenizer ([`~transformers.CLIPTokenizer`]):
A `CLIPTokenizer` to tokenize text.
unet ([`UNet2DConditionModel`]):
A `UNet2DConditionModel` to denoise the encoded image latents.
controlnet ([`ControlNetXSModel`]):
Provides additional conditioning to the `unet` during the denoising process.
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
model_cpu_offload_seq = "text_encoder->unet->vae>controlnet"
_optional_components = ["safety_checker", "feature_extractor"]
_exclude_from_cpu_offload = ["safety_checker"]
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
controlnet: ControlNetXSModel,
scheduler: KarrasDiffusionSchedulers,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPImageProcessor,
requires_safety_checker: bool = True,
):
super().__init__()
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
vae_compatible, cnxs_condition_downsample_factor, vae_downsample_factor = controlnet._check_if_vae_compatible(
vae
)
if not vae_compatible:
raise ValueError(
f"The downsampling factors of the VAE ({vae_downsample_factor}) and the conditioning part of ControlNetXS model {cnxs_condition_downsample_factor} need to be equal. Consider building the ControlNetXS model with different `conditioning_block_sizes`."
)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
controlnet=controlnet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor, do_convert_rgb=True)
self.control_image_processor = VaeImageProcessor(
vae_scale_factor=self.vae_scale_factor, do_convert_rgb=True, do_normalize=False
)
self.register_to_config(requires_safety_checker=requires_safety_checker)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_slicing
def enable_vae_slicing(self):
r"""
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
"""
self.vae.enable_slicing()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_slicing
def disable_vae_slicing(self):
r"""
Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_slicing()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_tiling
def enable_vae_tiling(self):
r"""
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow
processing larger images.
"""
self.vae.enable_tiling()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling
def disable_vae_tiling(self):
r"""
Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_tiling()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._encode_prompt
def _encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
**kwargs,
):
deprecation_message = "`_encode_prompt()` is deprecated and it will be removed in a future version. Use `encode_prompt()` instead. Also, be aware that the output format changed from a concatenated tensor to a tuple."
deprecate("_encode_prompt()", "1.0.0", deprecation_message, standard_warn=False)
prompt_embeds_tuple = self.encode_prompt(
prompt=prompt,
device=device,
num_images_per_prompt=num_images_per_prompt,
do_classifier_free_guidance=do_classifier_free_guidance,
negative_prompt=negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
lora_scale=lora_scale,
**kwargs,
)
# concatenate for backwards comp
prompt_embeds = torch.cat([prompt_embeds_tuple[1], prompt_embeds_tuple[0]])
return prompt_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.encode_prompt
def encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
clip_skip: Optional[int] = None,
):
r"""
Encodes the prompt into text encoder hidden states.
Args:
prompt (`str` or `List[str]`, *optional*):
prompt to be encoded
device: (`torch.device`):
torch device
num_images_per_prompt (`int`):
number of images that should be generated per prompt
do_classifier_free_guidance (`bool`):
whether to use classifier free guidance or not
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. If not defined, one has to pass
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is
less than `1`).
prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not
provided, text embeddings will be generated from `prompt` input argument.
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input
argument.
lora_scale (`float`, *optional*):
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(self.text_encoder, lora_scale)
else:
scale_lora_layers(self.text_encoder, lora_scale)
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
if prompt_embeds is None:
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
prompt = self.maybe_convert_prompt(prompt, self.tokenizer)
text_inputs = self.tokenizer(
prompt,
padding="max_length",
max_length=self.tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
untruncated_ids = self.tokenizer(prompt, padding="longest", return_tensors="pt").input_ids
if untruncated_ids.shape[-1] >= text_input_ids.shape[-1] and not torch.equal(
text_input_ids, untruncated_ids
):
removed_text = self.tokenizer.batch_decode(
untruncated_ids[:, self.tokenizer.model_max_length - 1 : -1]
)
logger.warning(
"The following part of your input was truncated because CLIP can only handle sequences up to"
f" {self.tokenizer.model_max_length} tokens: {removed_text}"
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = text_inputs.attention_mask.to(device)
else:
attention_mask = None
if clip_skip is None:
prompt_embeds = self.text_encoder(text_input_ids.to(device), attention_mask=attention_mask)
prompt_embeds = prompt_embeds[0]
else:
prompt_embeds = self.text_encoder(
text_input_ids.to(device), attention_mask=attention_mask, output_hidden_states=True
)
# Access the `hidden_states` first, that contains a tuple of
# all the hidden states from the encoder layers. Then index into
# the tuple to access the hidden states from the desired layer.
prompt_embeds = prompt_embeds[-1][-(clip_skip + 1)]
# We also need to apply the final LayerNorm here to not mess with the
# representations. The `last_hidden_states` that we typically use for
# obtaining the final prompt representations passes through the LayerNorm
# layer.
prompt_embeds = self.text_encoder.text_model.final_layer_norm(prompt_embeds)
if self.text_encoder is not None:
prompt_embeds_dtype = self.text_encoder.dtype
elif self.unet is not None:
prompt_embeds_dtype = self.unet.dtype
else:
prompt_embeds_dtype = prompt_embeds.dtype
prompt_embeds = prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
bs_embed, seq_len, _ = prompt_embeds.shape
# duplicate text embeddings for each generation per prompt, using mps friendly method
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1)
# get unconditional embeddings for classifier free guidance
if do_classifier_free_guidance and negative_prompt_embeds is None:
uncond_tokens: List[str]
if negative_prompt is None:
uncond_tokens = [""] * batch_size
elif prompt is not None and type(prompt) is not type(negative_prompt):
raise TypeError(
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !="
f" {type(prompt)}."
)
elif isinstance(negative_prompt, str):
uncond_tokens = [negative_prompt]
elif batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
else:
uncond_tokens = negative_prompt
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
uncond_tokens = self.maybe_convert_prompt(uncond_tokens, self.tokenizer)
max_length = prompt_embeds.shape[1]
uncond_input = self.tokenizer(
uncond_tokens,
padding="max_length",
max_length=max_length,
truncation=True,
return_tensors="pt",
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = uncond_input.attention_mask.to(device)
else:
attention_mask = None
negative_prompt_embeds = self.text_encoder(
uncond_input.input_ids.to(device),
attention_mask=attention_mask,
)
negative_prompt_embeds = negative_prompt_embeds[0]
if do_classifier_free_guidance:
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
seq_len = negative_prompt_embeds.shape[1]
negative_prompt_embeds = negative_prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
return prompt_embeds, negative_prompt_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker
def run_safety_checker(self, image, device, dtype):
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.decode_latents
def decode_latents(self, latents):
deprecation_message = "The decode_latents method is deprecated and will be removed in 1.0.0. Please use VaeImageProcessor.postprocess(...) instead"
deprecate("decode_latents", "1.0.0", deprecation_message, standard_warn=False)
latents = 1 / self.vae.config.scaling_factor * latents
image = self.vae.decode(latents, return_dict=False)[0]
image = (image / 2 + 0.5).clamp(0, 1)
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloat16
image = image.cpu().permute(0, 2, 3, 1).float().numpy()
return image
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs
def prepare_extra_step_kwargs(self, generator, eta):
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# check if the scheduler accepts generator
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
if accepts_generator:
extra_step_kwargs["generator"] = generator
return extra_step_kwargs
def check_inputs(
self,
prompt,
image,
callback_steps,
negative_prompt=None,
prompt_embeds=None,
negative_prompt_embeds=None,
controlnet_conditioning_scale=1.0,
control_guidance_start=0.0,
control_guidance_end=1.0,
):
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
if prompt is not None and prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to"
" only forward one of the two."
)
elif prompt is None and prompt_embeds is None:
raise ValueError(
"Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined."
)
elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if negative_prompt is not None and negative_prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:"
f" {negative_prompt_embeds}. Please make sure to only forward one of the two."
)
if prompt_embeds is not None and negative_prompt_embeds is not None:
if prompt_embeds.shape != negative_prompt_embeds.shape:
raise ValueError(
"`prompt_embeds` and `negative_prompt_embeds` must have the same shape when passed directly, but"
f" got: `prompt_embeds` {prompt_embeds.shape} != `negative_prompt_embeds`"
f" {negative_prompt_embeds.shape}."
)
# Check `image`
is_compiled = hasattr(F, "scaled_dot_product_attention") and isinstance(
self.controlnet, torch._dynamo.eval_frame.OptimizedModule
)
if (
isinstance(self.controlnet, ControlNetXSModel)
or is_compiled
and isinstance(self.controlnet._orig_mod, ControlNetXSModel)
):
self.check_image(image, prompt, prompt_embeds)
else:
assert False
# Check `controlnet_conditioning_scale`
if (
isinstance(self.controlnet, ControlNetXSModel)
or is_compiled
and isinstance(self.controlnet._orig_mod, ControlNetXSModel)
):
if not isinstance(controlnet_conditioning_scale, float):
raise TypeError("For single controlnet: `controlnet_conditioning_scale` must be type `float`.")
else:
assert False
start, end = control_guidance_start, control_guidance_end
if start >= end:
raise ValueError(
f"control guidance start: {start} cannot be larger or equal to control guidance end: {end}."
)
if start < 0.0:
raise ValueError(f"control guidance start: {start} can't be smaller than 0.")
if end > 1.0:
raise ValueError(f"control guidance end: {end} can't be larger than 1.0.")
def check_image(self, image, prompt, prompt_embeds):
image_is_pil = isinstance(image, PIL.Image.Image)
image_is_tensor = isinstance(image, torch.Tensor)
image_is_np = isinstance(image, np.ndarray)
image_is_pil_list = isinstance(image, list) and isinstance(image[0], PIL.Image.Image)
image_is_tensor_list = isinstance(image, list) and isinstance(image[0], torch.Tensor)
image_is_np_list = isinstance(image, list) and isinstance(image[0], np.ndarray)
if (
not image_is_pil
and not image_is_tensor
and not image_is_np
and not image_is_pil_list
and not image_is_tensor_list
and not image_is_np_list
):
raise TypeError(
f"image must be passed and be one of PIL image, numpy array, torch tensor, list of PIL images, list of numpy arrays or list of torch tensors, but is {type(image)}"
)
if image_is_pil:
image_batch_size = 1
else:
image_batch_size = len(image)
if prompt is not None and isinstance(prompt, str):
prompt_batch_size = 1
elif prompt is not None and isinstance(prompt, list):
prompt_batch_size = len(prompt)
elif prompt_embeds is not None:
prompt_batch_size = prompt_embeds.shape[0]
if image_batch_size != 1 and image_batch_size != prompt_batch_size:
raise ValueError(
f"If image batch size is not 1, image batch size must be same as prompt batch size. image batch size: {image_batch_size}, prompt batch size: {prompt_batch_size}"
)
def prepare_image(
self,
image,
width,
height,
batch_size,
num_images_per_prompt,
device,
dtype,
do_classifier_free_guidance=False,
):
image = self.control_image_processor.preprocess(image, height=height, width=width).to(dtype=torch.float32)
image_batch_size = image.shape[0]
if image_batch_size == 1:
repeat_by = batch_size
else:
# image batch size is the same as prompt batch size
repeat_by = num_images_per_prompt
image = image.repeat_interleave(repeat_by, dim=0)
image = image.to(device=device, dtype=dtype)
if do_classifier_free_guidance:
image = torch.cat([image] * 2)
return image
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
if latents is None:
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
else:
latents = latents.to(device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
return latents
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_freeu
def enable_freeu(self, s1: float, s2: float, b1: float, b2: float):
r"""Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497.
The suffixes after the scaling factors represent the stages where they are being applied.
Please refer to the [official repository](https://github.com/ChenyangSi/FreeU) for combinations of the values
that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL.
Args:
s1 (`float`):
Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to
mitigate "oversmoothing effect" in the enhanced denoising process.
s2 (`float`):
Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to
mitigate "oversmoothing effect" in the enhanced denoising process.
b1 (`float`): Scaling factor for stage 1 to amplify the contributions of backbone features.
b2 (`float`): Scaling factor for stage 2 to amplify the contributions of backbone features.
"""
if not hasattr(self, "unet"):
raise ValueError("The pipeline must have `unet` for using FreeU.")
self.unet.enable_freeu(s1=s1, s2=s2, b1=b1, b2=b2)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_freeu
def disable_freeu(self):
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
@torch.no_grad()
@replace_example_docstring(EXAMPLE_DOC_STRING)
def __call__(
self,
prompt: Union[str, List[str]] = None,
image: PipelineImageInput = None,
height: Optional[int] = None,
width: Optional[int] = None,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
latents: Optional[torch.FloatTensor] = None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: int = 1,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
controlnet_conditioning_scale: Union[float, List[float]] = 1.0,
control_guidance_start: float = 0.0,
control_guidance_end: float = 1.0,
clip_skip: Optional[int] = None,
):
r"""
The call function to the pipeline for generation.
Args:
prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide image generation. If not defined, you need to pass `prompt_embeds`.
image (`torch.FloatTensor`, `PIL.Image.Image`, `np.ndarray`, `List[torch.FloatTensor]`, `List[PIL.Image.Image]`, `List[np.ndarray]`,
`List[List[torch.FloatTensor]]`, `List[List[np.ndarray]]` or `List[List[PIL.Image.Image]]`):
The ControlNet input condition to provide guidance to the `unet` for generation. If the type is
specified as `torch.FloatTensor`, it is passed to ControlNet as is. `PIL.Image.Image` can also be
accepted as an image. The dimensions of the output image defaults to `image`'s dimensions. If height
and/or width are passed, `image` is resized accordingly. If multiple ControlNets are specified in
`init`, images must be passed as a list such that each element of the list can be correctly batched for
input to a single ControlNet.
height (`int`, *optional*, defaults to `self.unet.config.sample_size * self.vae_scale_factor`):
The height in pixels of the generated image.
width (`int`, *optional*, defaults to `self.unet.config.sample_size * self.vae_scale_factor`):
The width in pixels of the generated image.
num_inference_steps (`int`, *optional*, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
guidance_scale (`float`, *optional*, defaults to 7.5):
A higher guidance scale value encourages the model to generate images closely linked to the text
`prompt` at the expense of lower image quality. Guidance scale is enabled when `guidance_scale > 1`.
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass `negative_prompt_embeds` instead. Ignored when not using guidance (`guidance_scale < 1`).
num_images_per_prompt (`int`, *optional*, defaults to 1):
The number of images to generate per prompt.
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) from the [DDIM](https://arxiv.org/abs/2010.02502) paper. Only applies
to the [`~schedulers.DDIMScheduler`], and is ignored in other schedulers.
generator (`torch.Generator` or `List[torch.Generator]`, *optional*):
A [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make
generation deterministic.
latents (`torch.FloatTensor`, *optional*):
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random `generator`.
prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the `prompt` input argument.
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided, `negative_prompt_embeds` are generated from the `negative_prompt` input argument.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generated image. Choose between `PIL.Image` or `np.array`.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a
plain tuple.
callback (`Callable`, *optional*):
A function that calls every `callback_steps` steps during inference. The function is called with the
following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function is called. If not specified, the callback is called at
every step.
cross_attention_kwargs (`dict`, *optional*):
A kwargs dictionary that if specified is passed along to the [`AttentionProcessor`] as defined in
[`self.processor`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
controlnet_conditioning_scale (`float` or `List[float]`, *optional*, defaults to 1.0):
The outputs of the ControlNet are multiplied by `controlnet_conditioning_scale` before they are added
to the residual in the original `unet`. If multiple ControlNets are specified in `init`, you can set
the corresponding scale as a list.
control_guidance_start (`float` or `List[float]`, *optional*, defaults to 0.0):
The percentage of total steps at which the ControlNet starts applying.
control_guidance_end (`float` or `List[float]`, *optional*, defaults to 1.0):
The percentage of total steps at which the ControlNet stops applying.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
Examples:
Returns:
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`:
If `return_dict` is `True`, [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] is returned,
otherwise a `tuple` is returned where the first element is a list with the generated images and the
second element is a list of `bool`s indicating whether the corresponding generated image contains
"not-safe-for-work" (nsfw) content.
"""
controlnet = self.controlnet._orig_mod if is_compiled_module(self.controlnet) else self.controlnet
# 1. Check inputs. Raise error if not correct
self.check_inputs(
prompt,
image,
callback_steps,
negative_prompt,
prompt_embeds,
negative_prompt_embeds,
controlnet_conditioning_scale,
control_guidance_start,
control_guidance_end,
)
# 2. Define call parameters
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
device = self._execution_device
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# 3. Encode input prompt
text_encoder_lora_scale = (
cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None
)
prompt_embeds, negative_prompt_embeds = self.encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
lora_scale=text_encoder_lora_scale,
clip_skip=clip_skip,
)
# For classifier free guidance, we need to do two forward passes.
# Here we concatenate the unconditional and text embeddings into a single batch
# to avoid doing two forward passes
if do_classifier_free_guidance:
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds])
# 4. Prepare image
if isinstance(controlnet, ControlNetXSModel):
image = self.prepare_image(
image=image,
width=width,
height=height,
batch_size=batch_size * num_images_per_prompt,
num_images_per_prompt=num_images_per_prompt,
device=device,
dtype=controlnet.dtype,
do_classifier_free_guidance=do_classifier_free_guidance,
)
height, width = image.shape[-2:]
else:
assert False
# 5. Prepare timesteps
self.scheduler.set_timesteps(num_inference_steps, device=device)
timesteps = self.scheduler.timesteps
# 6. Prepare latent variables
num_channels_latents = self.unet.config.in_channels
latents = self.prepare_latents(
batch_size * num_images_per_prompt,
num_channels_latents,
height,
width,
prompt_embeds.dtype,
device,
generator,
latents,
)
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 8. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
is_unet_compiled = is_compiled_module(self.unet)
is_controlnet_compiled = is_compiled_module(self.controlnet)
is_torch_higher_equal_2_1 = is_torch_version(">=", "2.1")
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
# Relevant thread:
# https://dev-discuss.pytorch.org/t/cudagraphs-in-pytorch-2-0/1428
if (is_unet_compiled and is_controlnet_compiled) and is_torch_higher_equal_2_1:
torch._inductor.cudagraph_mark_step_begin()
# expand the latents if we are doing classifier free guidance
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
# predict the noise residual
dont_control = (
i / len(timesteps) < control_guidance_start or (i + 1) / len(timesteps) > control_guidance_end
)
if dont_control:
noise_pred = self.unet(
sample=latent_model_input,
timestep=t,
encoder_hidden_states=prompt_embeds,
cross_attention_kwargs=cross_attention_kwargs,
return_dict=True,
).sample
else:
noise_pred = self.controlnet(
base_model=self.unet,
sample=latent_model_input,
timestep=t,
encoder_hidden_states=prompt_embeds,
controlnet_cond=image,
conditioning_scale=controlnet_conditioning_scale,
cross_attention_kwargs=cross_attention_kwargs,
return_dict=True,
).sample
# perform guidance
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0]
# call the callback, if provided
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0):
progress_bar.update()
if callback is not None and i % callback_steps == 0:
step_idx = i // getattr(self.scheduler, "order", 1)
callback(step_idx, t, latents)
# If we do sequential model offloading, let's offload unet and controlnet
# manually for max memory savings
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.unet.to("cpu")
self.controlnet.to("cpu")
torch.cuda.empty_cache()
if not output_type == "latent":
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False, generator=generator)[
0
]
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
has_nsfw_concept = None
if has_nsfw_concept is None:
do_denormalize = [True] * image.shape[0]
else:
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept]
image = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
# Offload all models
self.maybe_free_model_hooks()
if not return_dict:
return (image, has_nsfw_concept)
return StableDiffusionPipelineOutput(images=image, nsfw_content_detected=has_nsfw_concept)
File diff suppressed because it is too large Load Diff
@@ -446,9 +446,8 @@ def convert_ldm_unet_checkpoint(
new_checkpoint["add_embedding.linear_2.bias"] = unet_state_dict["label_emb.0.2.bias"]
# Relevant to StableDiffusionUpscalePipeline
if "num_class_embeds" in config:
if (config["num_class_embeds"] is not None) and ("label_emb.weight" in unet_state_dict):
new_checkpoint["class_embedding.weight"] = unet_state_dict["label_emb.weight"]
if (config["num_class_embeds"] is not None) and ("label_emb.weight" in unet_state_dict):
new_checkpoint["class_embedding.weight"] = unet_state_dict["label_emb.weight"]
new_checkpoint["conv_in.weight"] = unet_state_dict["input_blocks.0.0.weight"]
new_checkpoint["conv_in.bias"] = unet_state_dict["input_blocks.0.0.bias"]
@@ -34,7 +34,6 @@ from ...loaders import (
from ...models import AutoencoderKL, ImageProjection, UNet2DConditionModel
from ...models.attention_processor import (
AttnProcessor2_0,
FusedAttnProcessor2_0,
LoRAAttnProcessor2_0,
LoRAXFormersAttnProcessor,
XFormersAttnProcessor,
@@ -682,6 +681,7 @@ class StableDiffusionXLPipeline(
add_time_ids = torch.tensor([add_time_ids], dtype=dtype)
return add_time_ids
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_upscale.StableDiffusionUpscalePipeline.upcast_vae
def upcast_vae(self):
dtype = self.vae.dtype
self.vae.to(dtype=torch.float32)
@@ -692,7 +692,6 @@ class StableDiffusionXLPipeline(
XFormersAttnProcessor,
LoRAXFormersAttnProcessor,
LoRAAttnProcessor2_0,
FusedAttnProcessor2_0,
),
)
# if xformers or torch_2_0 is used attention block does not need
@@ -730,65 +729,6 @@ class StableDiffusionXLPipeline(
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
def fuse_qkv_projections(self, unet: bool = True, vae: bool = True):
"""
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query,
key, value) are fused. For cross-attention modules, key and value projection matrices are fused.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
Args:
unet (`bool`, defaults to `True`): To apply fusion on the UNet.
vae (`bool`, defaults to `True`): To apply fusion on the VAE.
"""
self.fusing_unet = False
self.fusing_vae = False
if unet:
self.fusing_unet = True
self.unet.fuse_qkv_projections()
self.unet.set_attn_processor(FusedAttnProcessor2_0())
if vae:
if not isinstance(self.vae, AutoencoderKL):
raise ValueError("`fuse_qkv_projections()` is only supported for the VAE of type `AutoencoderKL`.")
self.fusing_vae = True
self.vae.fuse_qkv_projections()
self.vae.set_attn_processor(FusedAttnProcessor2_0())
def unfuse_qkv_projections(self, unet: bool = True, vae: bool = True):
"""Disable QKV projection fusion if enabled.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
Args:
unet (`bool`, defaults to `True`): To apply fusion on the UNet.
vae (`bool`, defaults to `True`): To apply fusion on the VAE.
"""
if unet:
if not self.fusing_unet:
logger.warning("The UNet was not initially fused for QKV projections. Doing nothing.")
else:
self.unet.unfuse_qkv_projections()
self.fusing_unet = False
if vae:
if not self.fusing_vae:
logger.warning("The VAE was not initially fused for QKV projections. Doing nothing.")
else:
self.vae.unfuse_qkv_projections()
self.fusing_vae = False
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
@@ -24,7 +24,6 @@ from ...loaders import FromSingleFileMixin, StableDiffusionXLLoraLoaderMixin, Te
from ...models import AutoencoderKL, UNet2DConditionModel
from ...models.attention_processor import (
AttnProcessor2_0,
FusedAttnProcessor2_0,
LoRAAttnProcessor2_0,
LoRAXFormersAttnProcessor,
XFormersAttnProcessor,
@@ -611,7 +610,6 @@ class StableDiffusionXLInstructPix2PixPipeline(
XFormersAttnProcessor,
LoRAXFormersAttnProcessor,
LoRAAttnProcessor2_0,
FusedAttnProcessor2_0,
),
)
# if xformers or torch_2_0 is used attention block does not need
@@ -10,10 +10,10 @@ from diffusers.utils import deprecate
from ...configuration_utils import ConfigMixin, register_to_config
from ...models import ModelMixin
from ...models.activations import get_activation
from ...models.attention import Attention
from ...models.attention_processor import (
ADDED_KV_ATTENTION_PROCESSORS,
CROSS_ATTENTION_PROCESSORS,
Attention,
AttentionProcessor,
AttnAddedKVProcessor,
AttnAddedKVProcessor2_0,
@@ -1000,42 +1000,6 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
if hasattr(upsample_block, k) or getattr(upsample_block, k, None) is not None:
setattr(upsample_block, k, None)
def fuse_qkv_projections(self):
"""
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query,
key, value) are fused. For cross-attention modules, key and value projection matrices are fused.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
self.original_attn_processors = None
for _, attn_processor in self.attn_processors.items():
if "Added" in str(attn_processor.__class__.__name__):
raise ValueError("`fuse_qkv_projections()` is not supported for models having added KV projections.")
self.original_attn_processors = self.attn_processors
for module in self.modules():
if isinstance(module, Attention):
module.fuse_projections(fuse=True)
def unfuse_qkv_projections(self):
"""Disables the fused QKV projection if enabled.
<Tip warning={true}>
This API is 🧪 experimental.
</Tip>
"""
if self.original_attn_processors is not None:
self.set_attn_processor(self.original_attn_processors)
def forward(
self,
sample: torch.FloatTensor,
@@ -734,16 +734,7 @@ class DEISMultistepScheduler(SchedulerMixin, ConfigMixin):
schedule_timesteps = self.timesteps.to(original_samples.device)
timesteps = timesteps.to(original_samples.device)
step_indices = []
for timestep in timesteps:
index_candidates = (schedule_timesteps == timestep).nonzero()
if len(index_candidates) == 0:
step_index = len(schedule_timesteps) - 1
elif len(index_candidates) > 1:
step_index = index_candidates[1].item()
else:
step_index = index_candidates[0].item()
step_indices.append(step_index)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < len(original_samples.shape):
@@ -896,16 +896,7 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
schedule_timesteps = self.timesteps.to(original_samples.device)
timesteps = timesteps.to(original_samples.device)
step_indices = []
for timestep in timesteps:
index_candidates = (schedule_timesteps == timestep).nonzero()
if len(index_candidates) == 0:
step_index = len(schedule_timesteps) - 1
elif len(index_candidates) > 1:
step_index = index_candidates[1].item()
else:
step_index = index_candidates[0].item()
step_indices.append(step_index)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < len(original_samples.shape):
@@ -891,16 +891,7 @@ class DPMSolverMultistepInverseScheduler(SchedulerMixin, ConfigMixin):
schedule_timesteps = self.timesteps.to(original_samples.device)
timesteps = timesteps.to(original_samples.device)
step_indices = []
for timestep in timesteps:
index_candidates = (schedule_timesteps == timestep).nonzero()
if len(index_candidates) == 0:
step_index = len(schedule_timesteps) - 1
elif len(index_candidates) > 1:
step_index = index_candidates[1].item()
else:
step_index = index_candidates[0].item()
step_indices.append(step_index)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < len(original_samples.shape):
@@ -897,16 +897,7 @@ class DPMSolverSinglestepScheduler(SchedulerMixin, ConfigMixin):
schedule_timesteps = self.timesteps.to(original_samples.device)
timesteps = timesteps.to(original_samples.device)
step_indices = []
for timestep in timesteps:
index_candidates = (schedule_timesteps == timestep).nonzero()
if len(index_candidates) == 0:
step_index = len(schedule_timesteps) - 1
elif len(index_candidates) > 1:
step_index = index_candidates[1].item()
else:
step_index = index_candidates[0].item()
step_indices.append(step_index)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < len(original_samples.shape):
@@ -92,43 +92,6 @@ def betas_for_alpha_bar(
return torch.tensor(betas, dtype=torch.float32)
# Copied from diffusers.schedulers.scheduling_ddim.rescale_zero_terminal_snr
def rescale_zero_terminal_snr(betas):
"""
Rescales betas to have zero terminal SNR Based on https://arxiv.org/pdf/2305.08891.pdf (Algorithm 1)
Args:
betas (`torch.FloatTensor`):
the betas that the scheduler is being initialized with.
Returns:
`torch.FloatTensor`: rescaled betas with zero terminal SNR
"""
# Convert betas to alphas_bar_sqrt
alphas = 1.0 - betas
alphas_cumprod = torch.cumprod(alphas, dim=0)
alphas_bar_sqrt = alphas_cumprod.sqrt()
# Store old values.
alphas_bar_sqrt_0 = alphas_bar_sqrt[0].clone()
alphas_bar_sqrt_T = alphas_bar_sqrt[-1].clone()
# Shift so the last timestep is zero.
alphas_bar_sqrt -= alphas_bar_sqrt_T
# Scale so the first timestep is back to the old value.
alphas_bar_sqrt *= alphas_bar_sqrt_0 / (alphas_bar_sqrt_0 - alphas_bar_sqrt_T)
# Convert alphas_bar_sqrt to betas
alphas_bar = alphas_bar_sqrt**2 # Revert sqrt
alphas = alphas_bar[1:] / alphas_bar[:-1] # Revert cumprod
alphas = torch.cat([alphas_bar[0:1], alphas])
betas = 1 - alphas
return betas
class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
"""
Euler scheduler.
@@ -165,10 +128,6 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
An offset added to the inference steps. You can use a combination of `offset=1` and
`set_alpha_to_one=False` to make the last step use step 0 for the previous alpha product like in Stable
Diffusion.
rescale_betas_zero_snr (`bool`, defaults to `False`):
Whether to rescale the betas to have zero terminal SNR. This enables the model to generate very bright and
dark samples instead of limiting it to samples with medium brightness. Loosely related to
[`--offset_noise`](https://github.com/huggingface/diffusers/blob/74fd735eb073eb1d774b1ab4154a0876eb82f055/examples/dreambooth/train_dreambooth.py#L506).
"""
_compatibles = [e.name for e in KarrasDiffusionSchedulers]
@@ -190,7 +149,6 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
timestep_spacing: str = "linspace",
timestep_type: str = "discrete", # can be "discrete" or "continuous"
steps_offset: int = 0,
rescale_betas_zero_snr: bool = False,
):
if trained_betas is not None:
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
@@ -205,17 +163,9 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
else:
raise NotImplementedError(f"{beta_schedule} does is not implemented for {self.__class__}")
if rescale_betas_zero_snr:
self.betas = rescale_zero_terminal_snr(self.betas)
self.alphas = 1.0 - self.betas
self.alphas_cumprod = torch.cumprod(self.alphas, dim=0)
if rescale_betas_zero_snr:
# Close to 0 without being 0 so first sigma is not inf
# FP16 smallest positive subnormal works well here
self.alphas_cumprod[-1] = 2**-24
sigmas = np.array(((1 - self.alphas_cumprod) / self.alphas_cumprod) ** 0.5)
timesteps = np.linspace(0, num_train_timesteps - 1, num_train_timesteps, dtype=float)[::-1].copy()
@@ -241,11 +191,10 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
@property
def init_noise_sigma(self):
# standard deviation of the initial noise distribution
max_sigma = max(self.sigmas) if isinstance(self.sigmas, list) else self.sigmas.max()
if self.config.timestep_spacing in ["linspace", "trailing"]:
return max_sigma
return self.sigmas.max()
return (max_sigma**2 + 1) ** 0.5
return (self.sigmas.max() ** 2 + 1) ** 0.5
@property
def step_index(self):
@@ -320,7 +269,7 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
if self.config.interpolation_type == "linear":
sigmas = np.interp(timesteps, np.arange(0, len(sigmas)), sigmas)
elif self.config.interpolation_type == "log_linear":
sigmas = torch.linspace(np.log(sigmas[-1]), np.log(sigmas[0]), num_inference_steps + 1).exp().numpy()
sigmas = torch.linspace(np.log(sigmas[-1]), np.log(sigmas[0]), num_inference_steps + 1).exp()
else:
raise ValueError(
f"{self.config.interpolation_type} is not implemented. Please specify interpolation_type to either"
@@ -470,9 +419,6 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
if self.step_index is None:
self._init_step_index(timestep)
# Upcast to avoid precision issues when computing prev_sample
sample = sample.to(torch.float32)
sigma = self.sigmas[self.step_index]
gamma = min(s_churn / (len(self.sigmas) - 1), 2**0.5 - 1) if s_tmin <= sigma <= s_tmax else 0.0
@@ -509,9 +455,6 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
prev_sample = sample + derivative * dt
# Cast sample back to model compatible dtype
prev_sample = prev_sample.to(model_output.dtype)
# upon completion increase step index by one
self._step_index += 1
@@ -828,16 +828,7 @@ class UniPCMultistepScheduler(SchedulerMixin, ConfigMixin):
schedule_timesteps = self.timesteps.to(original_samples.device)
timesteps = timesteps.to(original_samples.device)
step_indices = []
for timestep in timesteps:
index_candidates = (schedule_timesteps == timestep).nonzero()
if len(index_candidates) == 0:
step_index = len(schedule_timesteps) - 1
elif len(index_candidates) > 1:
step_index = index_candidates[1].item()
else:
step_index = index_candidates[0].item()
step_indices.append(step_index)
step_indices = [(schedule_timesteps == t).nonzero().item() for t in timesteps]
sigma = sigmas[step_indices].flatten()
while len(sigma.shape) < len(original_samples.shape):
-15
View File
@@ -92,21 +92,6 @@ class ControlNetModel(metaclass=DummyObject):
requires_backends(cls, ["torch"])
class ControlNetXSModel(metaclass=DummyObject):
_backends = ["torch"]
def __init__(self, *args, **kwargs):
requires_backends(self, ["torch"])
@classmethod
def from_config(cls, *args, **kwargs):
requires_backends(cls, ["torch"])
@classmethod
def from_pretrained(cls, *args, **kwargs):
requires_backends(cls, ["torch"])
class Kandinsky3UNet(metaclass=DummyObject):
_backends = ["torch"]
@@ -737,21 +737,6 @@ class StableDiffusionControlNetPipeline(metaclass=DummyObject):
requires_backends(cls, ["torch", "transformers"])
class StableDiffusionControlNetXSPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
def __init__(self, *args, **kwargs):
requires_backends(self, ["torch", "transformers"])
@classmethod
def from_config(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
@classmethod
def from_pretrained(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
class StableDiffusionDepth2ImgPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
@@ -1097,21 +1082,6 @@ class StableDiffusionXLControlNetPipeline(metaclass=DummyObject):
requires_backends(cls, ["torch", "transformers"])
class StableDiffusionXLControlNetXSPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
def __init__(self, *args, **kwargs):
requires_backends(self, ["torch", "transformers"])
@classmethod
def from_config(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
@classmethod
def from_pretrained(cls, *args, **kwargs):
requires_backends(cls, ["torch", "transformers"])
class StableDiffusionXLImg2ImgPipeline(metaclass=DummyObject):
_backends = ["torch", "transformers"]
+1 -1
View File
@@ -164,7 +164,7 @@ class PriorTransformerIntegrationTests(unittest.TestCase):
# clean up the VRAM after each test
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
@parameterized.expand(
[
@@ -869,7 +869,7 @@ class UNet2DConditionModelIntegrationTests(unittest.TestCase):
# clean up the VRAM after each test
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_latents(self, seed=0, shape=(4, 4, 64, 64), fp16=False):
dtype = torch.float16 if fp16 else torch.float32
+3 -3
View File
@@ -485,7 +485,7 @@ class AutoencoderTinyIntegrationTests(unittest.TestCase):
# clean up the VRAM after each test
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_file_format(self, seed, shape):
return f"gaussian_noise_s={seed}_shape={'_'.join([str(s) for s in shape])}.npy"
@@ -565,7 +565,7 @@ class AutoencoderKLIntegrationTests(unittest.TestCase):
# clean up the VRAM after each test
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_sd_image(self, seed=0, shape=(4, 3, 512, 512), fp16=False):
dtype = torch.float16 if fp16 else torch.float32
@@ -820,7 +820,7 @@ class AsymmetricAutoencoderKLIntegrationTests(unittest.TestCase):
# clean up the VRAM after each test
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_sd_image(self, seed=0, shape=(4, 3, 512, 512), fp16=False):
dtype = torch.float16 if fp16 else torch.float32
@@ -1,306 +0,0 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import gc
import traceback
import unittest
import numpy as np
import torch
from transformers import CLIPTextConfig, CLIPTextModel, CLIPTokenizer
from diffusers import (
AutoencoderKL,
ControlNetXSModel,
DDIMScheduler,
LCMScheduler,
StableDiffusionControlNetXSPipeline,
UNet2DConditionModel,
)
from diffusers.utils.import_utils import is_xformers_available
from diffusers.utils.testing_utils import (
enable_full_determinism,
load_image,
load_numpy,
require_python39_or_higher,
require_torch_2,
require_torch_gpu,
run_test_in_subprocess,
slow,
torch_device,
)
from diffusers.utils.torch_utils import randn_tensor
from ..pipeline_params import (
IMAGE_TO_IMAGE_IMAGE_PARAMS,
TEXT_TO_IMAGE_BATCH_PARAMS,
TEXT_TO_IMAGE_IMAGE_PARAMS,
TEXT_TO_IMAGE_PARAMS,
)
from ..test_pipelines_common import (
PipelineKarrasSchedulerTesterMixin,
PipelineLatentTesterMixin,
PipelineTesterMixin,
)
enable_full_determinism()
# Will be run via run_test_in_subprocess
def _test_stable_diffusion_compile(in_queue, out_queue, timeout):
error = None
try:
_ = in_queue.get(timeout=timeout)
controlnet = ControlNetXSModel.from_pretrained("UmerHA/ConrolNetXS-SD2.1-canny")
pipe = StableDiffusionControlNetXSPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", safety_checker=None, controlnet=controlnet
)
pipe.to("cuda")
pipe.set_progress_bar_config(disable=None)
pipe.unet.to(memory_format=torch.channels_last)
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipe.controlnet.to(memory_format=torch.channels_last)
pipe.controlnet = torch.compile(pipe.controlnet, mode="reduce-overhead", fullgraph=True)
generator = torch.Generator(device="cpu").manual_seed(0)
prompt = "bird"
image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny.png"
).resize((512, 512))
output = pipe(prompt, image, num_inference_steps=10, generator=generator, output_type="np")
image = output.images[0]
assert image.shape == (512, 512, 3)
expected_image = load_numpy(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny_out_full.npy"
)
expected_image = np.resize(expected_image, (512, 512, 3))
assert np.abs(expected_image - image).max() < 1.0
except Exception:
error = f"{traceback.format_exc()}"
results = {"error": error}
out_queue.put(results, timeout=timeout)
out_queue.join()
class ControlNetXSPipelineFastTests(
PipelineLatentTesterMixin, PipelineKarrasSchedulerTesterMixin, PipelineTesterMixin, unittest.TestCase
):
pipeline_class = StableDiffusionControlNetXSPipeline
params = TEXT_TO_IMAGE_PARAMS
batch_params = TEXT_TO_IMAGE_BATCH_PARAMS
image_params = IMAGE_TO_IMAGE_IMAGE_PARAMS
image_latents_params = TEXT_TO_IMAGE_IMAGE_PARAMS
def get_dummy_components(self, time_cond_proj_dim=None):
torch.manual_seed(0)
unet = UNet2DConditionModel(
block_out_channels=(4, 8),
layers_per_block=2,
sample_size=32,
in_channels=4,
out_channels=4,
down_block_types=("DownBlock2D", "CrossAttnDownBlock2D"),
up_block_types=("CrossAttnUpBlock2D", "UpBlock2D"),
cross_attention_dim=32,
norm_num_groups=1,
time_cond_proj_dim=time_cond_proj_dim,
)
torch.manual_seed(0)
controlnet = ControlNetXSModel.from_unet(
unet=unet,
time_embedding_mix=0.95,
learn_embedding=True,
size_ratio=0.5,
conditioning_embedding_out_channels=(16, 32),
num_attention_heads=2,
)
torch.manual_seed(0)
scheduler = DDIMScheduler(
beta_start=0.00085,
beta_end=0.012,
beta_schedule="scaled_linear",
clip_sample=False,
set_alpha_to_one=False,
)
torch.manual_seed(0)
vae = AutoencoderKL(
block_out_channels=[4, 8],
in_channels=3,
out_channels=3,
down_block_types=["DownEncoderBlock2D", "DownEncoderBlock2D"],
up_block_types=["UpDecoderBlock2D", "UpDecoderBlock2D"],
latent_channels=4,
norm_num_groups=2,
)
torch.manual_seed(0)
text_encoder_config = CLIPTextConfig(
bos_token_id=0,
eos_token_id=2,
hidden_size=32,
intermediate_size=37,
layer_norm_eps=1e-05,
num_attention_heads=4,
num_hidden_layers=5,
pad_token_id=1,
vocab_size=1000,
)
text_encoder = CLIPTextModel(text_encoder_config)
tokenizer = CLIPTokenizer.from_pretrained("hf-internal-testing/tiny-random-clip")
components = {
"unet": unet,
"controlnet": controlnet,
"scheduler": scheduler,
"vae": vae,
"text_encoder": text_encoder,
"tokenizer": tokenizer,
"safety_checker": None,
"feature_extractor": None,
}
return components
def get_dummy_inputs(self, device, seed=0):
if str(device).startswith("mps"):
generator = torch.manual_seed(seed)
else:
generator = torch.Generator(device=device).manual_seed(seed)
controlnet_embedder_scale_factor = 2
image = randn_tensor(
(1, 3, 32 * controlnet_embedder_scale_factor, 32 * controlnet_embedder_scale_factor),
generator=generator,
device=torch.device(device),
)
inputs = {
"prompt": "A painting of a squirrel eating a burger",
"generator": generator,
"num_inference_steps": 2,
"guidance_scale": 6.0,
"output_type": "numpy",
"image": image,
}
return inputs
def test_attention_slicing_forward_pass(self):
return self._test_attention_slicing_forward_pass(expected_max_diff=2e-3)
@unittest.skipIf(
torch_device != "cuda" or not is_xformers_available(),
reason="XFormers attention is only available with CUDA and `xformers` installed",
)
def test_xformers_attention_forwardGenerator_pass(self):
self._test_xformers_attention_forwardGenerator_pass(expected_max_diff=2e-3)
def test_inference_batch_single_identical(self):
self._test_inference_batch_single_identical(expected_max_diff=2e-3)
def test_controlnet_lcm(self):
device = "cpu" # ensure determinism for the device-dependent torch.Generator
components = self.get_dummy_components(time_cond_proj_dim=256)
sd_pipe = StableDiffusionControlNetXSPipeline(**components)
sd_pipe.scheduler = LCMScheduler.from_config(sd_pipe.scheduler.config)
sd_pipe = sd_pipe.to(torch_device)
sd_pipe.set_progress_bar_config(disable=None)
inputs = self.get_dummy_inputs(device)
output = sd_pipe(**inputs)
image = output.images
image_slice = image[0, -3:, -3:, -1]
assert image.shape == (1, 64, 64, 3)
expected_slice = np.array(
[0.52700454, 0.3930534, 0.25509018, 0.7132304, 0.53696585, 0.46568912, 0.7095368, 0.7059624, 0.4744786]
)
assert np.abs(image_slice.flatten() - expected_slice).max() < 1e-2
@slow
@require_torch_gpu
class ControlNetXSPipelineSlowTests(unittest.TestCase):
def tearDown(self):
super().tearDown()
gc.collect()
torch.cuda.empty_cache()
def test_canny(self):
controlnet = ControlNetXSModel.from_pretrained("UmerHA/ConrolNetXS-SD2.1-canny")
pipe = StableDiffusionControlNetXSPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", safety_checker=None, controlnet=controlnet
)
pipe.enable_model_cpu_offload()
pipe.set_progress_bar_config(disable=None)
generator = torch.Generator(device="cpu").manual_seed(0)
prompt = "bird"
image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny.png"
)
output = pipe(prompt, image, generator=generator, output_type="np", num_inference_steps=3)
image = output.images[0]
assert image.shape == (768, 512, 3)
original_image = image[-3:, -3:, -1].flatten()
expected_image = np.array([0.1274, 0.1401, 0.147, 0.1185, 0.1555, 0.1492, 0.1565, 0.1474, 0.1701])
assert np.allclose(original_image, expected_image, atol=1e-04)
def test_depth(self):
controlnet = ControlNetXSModel.from_pretrained("UmerHA/ConrolNetXS-SD2.1-depth")
pipe = StableDiffusionControlNetXSPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", safety_checker=None, controlnet=controlnet
)
pipe.enable_model_cpu_offload()
pipe.set_progress_bar_config(disable=None)
generator = torch.Generator(device="cpu").manual_seed(0)
prompt = "Stormtrooper's lecture"
image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/stormtrooper_depth.png"
)
output = pipe(prompt, image, generator=generator, output_type="np", num_inference_steps=3)
image = output.images[0]
assert image.shape == (512, 512, 3)
original_image = image[-3:, -3:, -1].flatten()
expected_image = np.array([0.1098, 0.1025, 0.1211, 0.1129, 0.1165, 0.1262, 0.1185, 0.1261, 0.1703])
assert np.allclose(original_image, expected_image, atol=1e-04)
@require_python39_or_higher
@require_torch_2
def test_stable_diffusion_compile(self):
run_test_in_subprocess(test_case=self, target_func=_test_stable_diffusion_compile, inputs=None)
@@ -1,362 +0,0 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import gc
import unittest
import numpy as np
import torch
from transformers import CLIPTextConfig, CLIPTextModel, CLIPTextModelWithProjection, CLIPTokenizer
from diffusers import (
AutoencoderKL,
ControlNetXSModel,
EulerDiscreteScheduler,
StableDiffusionXLControlNetXSPipeline,
UNet2DConditionModel,
)
from diffusers.utils.import_utils import is_xformers_available
from diffusers.utils.testing_utils import enable_full_determinism, load_image, require_torch_gpu, slow, torch_device
from diffusers.utils.torch_utils import randn_tensor
from ..pipeline_params import (
IMAGE_TO_IMAGE_IMAGE_PARAMS,
TEXT_TO_IMAGE_BATCH_PARAMS,
TEXT_TO_IMAGE_IMAGE_PARAMS,
TEXT_TO_IMAGE_PARAMS,
)
from ..test_pipelines_common import (
PipelineKarrasSchedulerTesterMixin,
PipelineLatentTesterMixin,
PipelineTesterMixin,
SDXLOptionalComponentsTesterMixin,
)
enable_full_determinism()
class StableDiffusionXLControlNetXSPipelineFastTests(
PipelineLatentTesterMixin,
PipelineKarrasSchedulerTesterMixin,
PipelineTesterMixin,
SDXLOptionalComponentsTesterMixin,
unittest.TestCase,
):
pipeline_class = StableDiffusionXLControlNetXSPipeline
params = TEXT_TO_IMAGE_PARAMS
batch_params = TEXT_TO_IMAGE_BATCH_PARAMS
image_params = IMAGE_TO_IMAGE_IMAGE_PARAMS
image_latents_params = TEXT_TO_IMAGE_IMAGE_PARAMS
def get_dummy_components(self):
torch.manual_seed(0)
unet = UNet2DConditionModel(
block_out_channels=(32, 64),
layers_per_block=2,
sample_size=32,
in_channels=4,
out_channels=4,
down_block_types=("DownBlock2D", "CrossAttnDownBlock2D"),
up_block_types=("CrossAttnUpBlock2D", "UpBlock2D"),
# SD2-specific config below
attention_head_dim=(2, 4),
use_linear_projection=True,
addition_embed_type="text_time",
addition_time_embed_dim=8,
transformer_layers_per_block=(1, 2),
projection_class_embeddings_input_dim=80, # 6 * 8 + 32
cross_attention_dim=64,
)
torch.manual_seed(0)
controlnet = ControlNetXSModel.from_unet(
unet,
time_embedding_mix=0.95,
learn_embedding=True,
size_ratio=0.5,
conditioning_embedding_out_channels=(16, 32),
)
torch.manual_seed(0)
scheduler = EulerDiscreteScheduler(
beta_start=0.00085,
beta_end=0.012,
steps_offset=1,
beta_schedule="scaled_linear",
timestep_spacing="leading",
)
torch.manual_seed(0)
vae = AutoencoderKL(
block_out_channels=[32, 64],
in_channels=3,
out_channels=3,
down_block_types=["DownEncoderBlock2D", "DownEncoderBlock2D"],
up_block_types=["UpDecoderBlock2D", "UpDecoderBlock2D"],
latent_channels=4,
)
torch.manual_seed(0)
text_encoder_config = CLIPTextConfig(
bos_token_id=0,
eos_token_id=2,
hidden_size=32,
intermediate_size=37,
layer_norm_eps=1e-05,
num_attention_heads=4,
num_hidden_layers=5,
pad_token_id=1,
vocab_size=1000,
# SD2-specific config below
hidden_act="gelu",
projection_dim=32,
)
text_encoder = CLIPTextModel(text_encoder_config)
tokenizer = CLIPTokenizer.from_pretrained("hf-internal-testing/tiny-random-clip")
text_encoder_2 = CLIPTextModelWithProjection(text_encoder_config)
tokenizer_2 = CLIPTokenizer.from_pretrained("hf-internal-testing/tiny-random-clip")
components = {
"unet": unet,
"controlnet": controlnet,
"scheduler": scheduler,
"vae": vae,
"text_encoder": text_encoder,
"tokenizer": tokenizer,
"text_encoder_2": text_encoder_2,
"tokenizer_2": tokenizer_2,
}
return components
# copied from test_controlnet_sdxl.py
def get_dummy_inputs(self, device, seed=0):
if str(device).startswith("mps"):
generator = torch.manual_seed(seed)
else:
generator = torch.Generator(device=device).manual_seed(seed)
controlnet_embedder_scale_factor = 2
image = randn_tensor(
(1, 3, 32 * controlnet_embedder_scale_factor, 32 * controlnet_embedder_scale_factor),
generator=generator,
device=torch.device(device),
)
inputs = {
"prompt": "A painting of a squirrel eating a burger",
"generator": generator,
"num_inference_steps": 2,
"guidance_scale": 6.0,
"output_type": "np",
"image": image,
}
return inputs
# copied from test_controlnet_sdxl.py
def test_attention_slicing_forward_pass(self):
return self._test_attention_slicing_forward_pass(expected_max_diff=2e-3)
# copied from test_controlnet_sdxl.py
@unittest.skipIf(
torch_device != "cuda" or not is_xformers_available(),
reason="XFormers attention is only available with CUDA and `xformers` installed",
)
def test_xformers_attention_forwardGenerator_pass(self):
self._test_xformers_attention_forwardGenerator_pass(expected_max_diff=2e-3)
# copied from test_controlnet_sdxl.py
def test_inference_batch_single_identical(self):
self._test_inference_batch_single_identical(expected_max_diff=2e-3)
# copied from test_controlnet_sdxl.py
def test_save_load_optional_components(self):
self._test_save_load_optional_components()
# copied from test_controlnet_sdxl.py
@require_torch_gpu
def test_stable_diffusion_xl_offloads(self):
pipes = []
components = self.get_dummy_components()
sd_pipe = self.pipeline_class(**components).to(torch_device)
pipes.append(sd_pipe)
components = self.get_dummy_components()
sd_pipe = self.pipeline_class(**components)
sd_pipe.enable_model_cpu_offload()
pipes.append(sd_pipe)
components = self.get_dummy_components()
sd_pipe = self.pipeline_class(**components)
sd_pipe.enable_sequential_cpu_offload()
pipes.append(sd_pipe)
image_slices = []
for pipe in pipes:
pipe.unet.set_default_attn_processor()
inputs = self.get_dummy_inputs(torch_device)
image = pipe(**inputs).images
image_slices.append(image[0, -3:, -3:, -1].flatten())
assert np.abs(image_slices[0] - image_slices[1]).max() < 1e-3
assert np.abs(image_slices[0] - image_slices[2]).max() < 1e-3
# copied from test_controlnet_sdxl.py
def test_stable_diffusion_xl_multi_prompts(self):
components = self.get_dummy_components()
sd_pipe = self.pipeline_class(**components).to(torch_device)
# forward with single prompt
inputs = self.get_dummy_inputs(torch_device)
output = sd_pipe(**inputs)
image_slice_1 = output.images[0, -3:, -3:, -1]
# forward with same prompt duplicated
inputs = self.get_dummy_inputs(torch_device)
inputs["prompt_2"] = inputs["prompt"]
output = sd_pipe(**inputs)
image_slice_2 = output.images[0, -3:, -3:, -1]
# ensure the results are equal
assert np.abs(image_slice_1.flatten() - image_slice_2.flatten()).max() < 1e-4
# forward with different prompt
inputs = self.get_dummy_inputs(torch_device)
inputs["prompt_2"] = "different prompt"
output = sd_pipe(**inputs)
image_slice_3 = output.images[0, -3:, -3:, -1]
# ensure the results are not equal
assert np.abs(image_slice_1.flatten() - image_slice_3.flatten()).max() > 1e-4
# manually set a negative_prompt
inputs = self.get_dummy_inputs(torch_device)
inputs["negative_prompt"] = "negative prompt"
output = sd_pipe(**inputs)
image_slice_1 = output.images[0, -3:, -3:, -1]
# forward with same negative_prompt duplicated
inputs = self.get_dummy_inputs(torch_device)
inputs["negative_prompt"] = "negative prompt"
inputs["negative_prompt_2"] = inputs["negative_prompt"]
output = sd_pipe(**inputs)
image_slice_2 = output.images[0, -3:, -3:, -1]
# ensure the results are equal
assert np.abs(image_slice_1.flatten() - image_slice_2.flatten()).max() < 1e-4
# forward with different negative_prompt
inputs = self.get_dummy_inputs(torch_device)
inputs["negative_prompt"] = "negative prompt"
inputs["negative_prompt_2"] = "different negative prompt"
output = sd_pipe(**inputs)
image_slice_3 = output.images[0, -3:, -3:, -1]
# ensure the results are not equal
assert np.abs(image_slice_1.flatten() - image_slice_3.flatten()).max() > 1e-4
# copied from test_stable_diffusion_xl.py
def test_stable_diffusion_xl_prompt_embeds(self):
components = self.get_dummy_components()
sd_pipe = self.pipeline_class(**components)
sd_pipe = sd_pipe.to(torch_device)
sd_pipe = sd_pipe.to(torch_device)
sd_pipe.set_progress_bar_config(disable=None)
# forward without prompt embeds
inputs = self.get_dummy_inputs(torch_device)
inputs["prompt"] = 2 * [inputs["prompt"]]
inputs["num_images_per_prompt"] = 2
output = sd_pipe(**inputs)
image_slice_1 = output.images[0, -3:, -3:, -1]
# forward with prompt embeds
inputs = self.get_dummy_inputs(torch_device)
prompt = 2 * [inputs.pop("prompt")]
(
prompt_embeds,
negative_prompt_embeds,
pooled_prompt_embeds,
negative_pooled_prompt_embeds,
) = sd_pipe.encode_prompt(prompt)
output = sd_pipe(
**inputs,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
pooled_prompt_embeds=pooled_prompt_embeds,
negative_pooled_prompt_embeds=negative_pooled_prompt_embeds,
)
image_slice_2 = output.images[0, -3:, -3:, -1]
# make sure that it's equal
assert np.abs(image_slice_1.flatten() - image_slice_2.flatten()).max() < 1.1e-4
@slow
@require_torch_gpu
class ControlNetSDXLPipelineXSSlowTests(unittest.TestCase):
def tearDown(self):
super().tearDown()
gc.collect()
torch.cuda.empty_cache()
def test_canny(self):
controlnet = ControlNetXSModel.from_pretrained("UmerHA/ConrolNetXS-SDXL-canny")
pipe = StableDiffusionXLControlNetXSPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", controlnet=controlnet
)
pipe.enable_sequential_cpu_offload()
pipe.set_progress_bar_config(disable=None)
generator = torch.Generator(device="cpu").manual_seed(0)
prompt = "bird"
image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny.png"
)
images = pipe(prompt, image=image, generator=generator, output_type="np", num_inference_steps=3).images
assert images[0].shape == (768, 512, 3)
original_image = images[0, -3:, -3:, -1].flatten()
expected_image = np.array([0.4359, 0.4335, 0.4609, 0.4515, 0.4669, 0.4494, 0.452, 0.4493, 0.4382])
assert np.allclose(original_image, expected_image, atol=1e-04)
def test_depth(self):
controlnet = ControlNetXSModel.from_pretrained("UmerHA/ConrolNetXS-SDXL-depth")
pipe = StableDiffusionXLControlNetXSPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", controlnet=controlnet
)
pipe.enable_sequential_cpu_offload()
pipe.set_progress_bar_config(disable=None)
generator = torch.Generator(device="cpu").manual_seed(0)
prompt = "Stormtrooper's lecture"
image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/stormtrooper_depth.png"
)
images = pipe(prompt, image=image, generator=generator, output_type="np", num_inference_steps=3).images
assert images[0].shape == (512, 512, 3)
original_image = images[0, -3:, -3:, -1].flatten()
expected_image = np.array([0.4411, 0.3617, 0.2654, 0.266, 0.3449, 0.3898, 0.3745, 0.353, 0.326])
assert np.allclose(original_image, expected_image, atol=1e-04)
@@ -182,25 +182,6 @@ class IPAdapterSDIntegrationTests(IPAdapterNightlyTestsMixin):
assert np.allclose(image_slice, expected_slice, atol=1e-4, rtol=1e-4)
def test_text_to_image_full_face(self):
image_encoder = self.get_image_encoder(repo_id="h94/IP-Adapter", subfolder="models/image_encoder")
pipeline = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", image_encoder=image_encoder, safety_checker=None, torch_dtype=self.dtype
)
pipeline.to(torch_device)
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter-full-face_sd15.bin")
pipeline.set_ip_adapter_scale(0.7)
inputs = self.get_dummy_inputs()
images = pipeline(**inputs).images
image_slice = images[0, :3, :3, -1].flatten()
expected_slice = np.array(
[0.1706543, 0.1303711, 0.12573242, 0.21777344, 0.14550781, 0.14038086, 0.40820312, 0.41455078, 0.42529297]
)
assert np.allclose(image_slice, expected_slice, atol=1e-4, rtol=1e-4)
@slow
@require_torch_gpu
@@ -310,7 +310,7 @@ class StableDiffusion2PipelineSlowTests(unittest.TestCase):
def tearDown(self):
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_inputs(self, device, generator_device="cpu", dtype=torch.float32, seed=0):
_generator_device = "cpu" if not generator_device.startswith("cuda") else "cuda"
@@ -531,7 +531,7 @@ class StableDiffusion2PipelineNightlyTests(unittest.TestCase):
def tearDown(self):
super().tearDown()
gc.collect()
backend_empty_cache(torch_device)
backend_empty_cache()
def get_inputs(self, device, generator_device="cpu", dtype=torch.float32, seed=0):
_generator_device = "cpu" if not generator_device.startswith("cuda") else "cuda"
@@ -938,37 +938,6 @@ class StableDiffusionXLPipelineFastTests(
assert np.abs(image_slices[0] - image_slices[1]).max() < 1e-3
def test_stable_diffusion_xl_with_fused_qkv_projections(self):
device = "cpu" # ensure determinism for the device-dependent torch.Generator
components = self.get_dummy_components()
sd_pipe = StableDiffusionXLPipeline(**components)
sd_pipe = sd_pipe.to(device)
sd_pipe.set_progress_bar_config(disable=None)
inputs = self.get_dummy_inputs(device)
image = sd_pipe(**inputs).images
original_image_slice = image[0, -3:, -3:, -1]
sd_pipe.fuse_qkv_projections()
inputs = self.get_dummy_inputs(device)
image = sd_pipe(**inputs).images
image_slice_fused = image[0, -3:, -3:, -1]
sd_pipe.unfuse_qkv_projections()
inputs = self.get_dummy_inputs(device)
image = sd_pipe(**inputs).images
image_slice_disabled = image[0, -3:, -3:, -1]
assert np.allclose(
original_image_slice, image_slice_fused, atol=1e-2, rtol=1e-2
), "Fusion of QKV projections shouldn't affect the outputs."
assert np.allclose(
image_slice_fused, image_slice_disabled, atol=1e-2, rtol=1e-2
), "Outputs, with QKV projection fusion enabled, shouldn't change when fused QKV projections are disabled."
assert np.allclose(
original_image_slice, image_slice_disabled, atol=1e-2, rtol=1e-2
), "Original outputs should match when fused QKV projections are disabled."
@slow
class StableDiffusionXLPipelineIntegrationTests(unittest.TestCase):
-4
View File
@@ -45,10 +45,6 @@ class EulerDiscreteSchedulerTest(SchedulerCommonTest):
def test_karras_sigmas(self):
self.check_over_configs(use_karras_sigmas=True, sigma_min=0.02, sigma_max=700.0)
def test_rescale_betas_zero_snr(self):
for rescale_betas_zero_snr in [True, False]:
self.check_over_configs(rescale_betas_zero_snr=rescale_betas_zero_snr)
def test_full_loop_no_noise(self):
scheduler_class = self.scheduler_classes[0]
scheduler_config = self.get_scheduler_config()