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110 Commits

Author SHA1 Message Date
Sayak Paul b08fc2d9d5 Merge branch 'main' into examples-test-fix 2023-12-14 21:49:52 +05:30
Monohydroxides c46711e895 [Community] Add SDE Drag pipeline (#6105)
* Add community pipeline: sde_drag.py

* Update README.md

* Update README.md

Update example code and visual example

* Update sde_drag.py

Update code example.
2023-12-14 20:47:20 +05:30
Dhruv Nair a031abdc89 add peft to training deps 2023-12-12 14:07:45 +00:00
Sayak Paul 1d686bac81 [feat: Benchmarking Workflow] add stuff for a benchmarking workflow (#5839)
* add poc for benchmarking workflow.

* import

* fix argument

* fix: argument

* fix: path

* fix

* fix

* path

* output csv files.

* workflow cleanup

* append token

* add utility to push to hf dataset

* fix: kw arg

* better reporting

* fix: headers

* better formatting of the numbers.

* better type annotation

* fix: formatting

* moentarily disable check

* push results.

* remove disable check

* introduce base classes.

* img2img class

* add inpainting pipeline

* intoduce base benchmark class.

* add img2img and inpainting

* feat: utility to compare changes

* fix

* fix import

* add args

* basepath

* better exception handling

* better path handling

* fix

* fix

* remove

* ifx

* fix

* add: support for controlnet.

* image_url -> url

* move images to huggingface hub

* correct urls.

* root_ckpt

* flush before benchmarking

* don't install accelerate from source

* add runner

* simplify Diffusers Benchmarking step

* change runner

* fix: subprocess call.

* filter percentage values

* fix controlnet benchmark

* add t2i adapters.

* fix filter columns

* fix t2i adapter benchmark

* fix init.

* fix

* remove safetensors flag

* fix args print

* fix

* feat: run_command

* add adapter resolution mapping

* benchmark t2i adapter fix.

* fix adapter input

* fix

* convert to L.

* add flush() add appropriate places

* better filtering

* okay

* get env for torch

* convert to float

* fix

* filter out nans.

* better coment

* sdxl

* sdxl for other benchmarks.

* fix: condition

* fix: condition for inpainting

* fix: mapping for resolution

* fix

* include kandinsky and wuerstchen

* fix: Wuerstchen

* Empty-Commit

* [Community] AnimateDiff + Controlnet Pipeline (#5928)

* begin work on animatediff + controlnet pipeline

* complete todos, uncomment multicontrolnet, input checks

Co-Authored-By: EdoardoBotta <botta.edoardo@gmail.com>

* update

Co-Authored-By: EdoardoBotta <botta.edoardo@gmail.com>

* add example

* update community README

* Update examples/community/README.md

---------

Co-authored-by: EdoardoBotta <botta.edoardo@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* EulerDiscreteScheduler add `rescale_betas_zero_snr` (#6024)

* EulerDiscreteScheduler add `rescale_betas_zero_snr`

* Revert "[Community] AnimateDiff + Controlnet Pipeline (#5928)"

This reverts commit 821726d7c0.

* Revert "EulerDiscreteScheduler add `rescale_betas_zero_snr` (#6024)"

This reverts commit 3dc2362b5a.

* add SDXL turbo

* add lcm lora to the mix as well.

* fix

* increase steps to 2 when running turbo i2i

* debug

* debug

* debug

* fix for good

* fix and isolate better

* fuse lora so that torch compile works with peft

* fix: LCMLoRA

* better identification for LCM

* change to cron job

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: Aryan V S <contact.aryanvs@gmail.com>
Co-authored-by: EdoardoBotta <botta.edoardo@gmail.com>
Co-authored-by: Beinsezii <39478211+Beinsezii@users.noreply.github.com>
2023-12-12 11:03:34 +05:30
M. Tolga Cangöz 0a401b95b7 [Docs] Fix typos (#6122)
Fix typos and trim trailing whitespaces
2023-12-11 10:55:28 -08:00
Edward Li 664e931bcb Correct type annotation for VaeImageProcessor.numpy_to_pil (#6111)
From `(np.ndarray) -> PIL.Image.Image` to `(np.ndarray) -> List[PIL.Image.Image]`.
2023-12-11 15:22:04 +05:30
Aryan V S 88bdd97ccd IP adapter support for most pipelines (#5900)
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_upscale.py

* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_attend_and_excite.py

* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py

* update tests

* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_panorama.py

* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_sag.py

* support ip-adapter in src/diffusers/pipelines/stable_diffusion_safe/pipeline_stable_diffusion_safe.py

* support ip-adapter in src/diffusers/pipelines/latent_consistency_models/pipeline_latent_consistency_text2img.py

* support ip-adapter in src/diffusers/pipelines/latent_consistency_models/pipeline_latent_consistency_img2img.py

* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_ldm3d.py

* revert changes to sd_attend_and_excite and sd_upscale

* make style

* fix broken tests

* update ip-adapter implementation to latest

* apply suggestions from review

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-10 21:19:14 +05:30
Charchit Sharma 08b453e382 IP-Adapter for StableDiffusionControlNetImg2ImgPipeline (#5901)
* adapter for StableDiffusionControlNetImg2ImgPipeline

* fix-copies

* fix-copies

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-09 11:02:55 +05:30
apolinário 2a111bc9fe [Advanced Training Script] Fix pipe example (#6106) 2023-12-08 15:56:35 +01:00
apolinário 16e6997f0d [Advanced Diffusion Script] Add Widget default text (#6100)
add widget
2023-12-08 12:45:27 +01:00
YiYi Xu 3b9b98656e Fix a bug in add_noise function (#6085)
* fix

* copies

---------

Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-12-07 11:35:28 -10:00
Fabio Rigano b65928b556 Add support for IPAdapterFull (#5911)
* Add support for IPAdapterFull


Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-07 06:40:39 -10:00
Beinsezii 6bf1ca2c79 EulerDiscreteScheduler add rescale_betas_zero_snr (#6024)
* EulerDiscreteScheduler add `rescale_betas_zero_snr`
2023-12-06 21:51:04 -10:00
Aryan V S 978dec9014 [Community] AnimateDiff + Controlnet Pipeline (#5928)
* begin work on animatediff + controlnet pipeline

* complete todos, uncomment multicontrolnet, input checks

Co-Authored-By: EdoardoBotta <botta.edoardo@gmail.com>

* update

Co-Authored-By: EdoardoBotta <botta.edoardo@gmail.com>

* add example

* update community README

* Update examples/community/README.md

---------

Co-authored-by: EdoardoBotta <botta.edoardo@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-06 21:01:41 -10:00
Dhruv Nair 79a7ab92d1 Fix clearing backend cache from device agnostic testing (#6075)
update
2023-12-07 11:18:31 +05:30
Younes Belkada c2717317f0 [PEFT] Adapt example scripts to use PEFT (#5388)
* adapt example scripts to use PEFT

* Update examples/text_to_image/train_text_to_image_lora.py

* fix

* add for SDXL

* oops

* make sure to install peft

* fix

* fix

* fix dreambooth and lora

* more fixes

* add peft to requirements.txt

* fix

* final fix

* add peft version in requirements

* remove comment

* change variable names

* add few lines in readme

* add to reqs

* style

* fix issues

* fix lora dreambooth xl tests

* init_lora_weights to gaussian and add out proj where missing

* ammend requirements.

* ammend requirements.txt

* add correct peft versions

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-07 09:39:29 +05:30
Ian bf7f9b49a2 Fix typing inconsistency in Euler discrete scheduler (#6052) 2023-12-06 23:45:16 +01:00
UmerHA e192ae08d3 Add ControlNet-XS support (#5827)
* Check in 23-10-05

* check-in 23-10-06

* check-in 23-10-07 2pm

* check-in 23-10-08

* check-in 231009T1200

* check-in 230109

* checkin 231010

* init + forward run

* checkin

* checkin

* ControlNetXSModel is now saveable+loadable

* Forward works

* checkin

* Pipeline works with `no_control=True`

* checkin

* debug: save intermediate outputs of resnet

* checkin

* Understood time error + fixed connection error

* checkin

* checkin 231106T1600

* turned off detailled debug prints

* time debug logs

* small fix

* Separated control_scale for connections/time

* simplified debug logging

* Full denoising works with control scale = 0

* aligned logs

* Added control_attention_head_dim param

* Passing n_heads instead of dim_head into ctrl unet

* Fixed ctrl midblock bug

* Cleanup

* Fixed time dtype bug

* checkin

* 1. from_unet, 2. base passed, 3. all unet params

* checkin

* Finished docstrings

* cleanup

* make style

* checkin

* more tests pass

* Fixed tests

* removed debug logs

* make style + quality

* make fix-copies

* fixed documentation

* added cnxs to doc toc

* added control start/end param

* Update controlnetxs_sdxl.md

* tried to fix copies..

* Fixed norm_num_groups in from_unet

* added sdxl-depth test

* created SD2.1 controlnet-xs pipeline

* re-added debug logs

* Adjusting group norm ; readded logs

* Added debug log statements

* removed debug logs ; started tests for sd2.1

* updated sd21 tests

* fixed tests

* fixed tests

* slightly increased error tolerance for 1 test

* make style & quality

* Added docs for CNXS-SD

* make fix-copies

* Fixed sd compile test ; fixed gradient ckpointing

* vae downs = cnxs conditioning downs; removed guess

* make style & quality

* Fixed tests

* fixed test

* Incorporated review feedback

* simplified control model surgery

* fixed tests & make style / quality

* Updated docs; deleted pip & cursor files

* Rolled back minimal change to resnet

* Update resnet.py

* Update resnet.py

* Update src/diffusers/models/controlnetxs.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update src/diffusers/models/controlnetxs.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Incorporated review feedback

* Update docs/source/en/api/pipelines/controlnetxs_sdxl.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/controlnetxs.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/controlnetxs.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/controlnetxs.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update src/diffusers/models/controlnetxs.py

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update src/diffusers/models/controlnetxs.py

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update src/diffusers/pipelines/controlnet_xs/pipeline_controlnet_xs.py

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/controlnetxs.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update src/diffusers/pipelines/controlnet_xs/pipeline_controlnet_xs_sd_xl.py

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Incorporated doc feedback

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2023-12-06 23:33:47 +01:00
Steven Liu 87a09d66f3 [docs] SDXL Turbo (#6065)
api docs
2023-12-06 14:33:14 -08:00
Lucain 75ada25048 Harmonize HF environment variables + deprecate use_auth_token (#6066)
* Harmonize HF environment variables + deprecate use_auth_token

* fix import

* fix
2023-12-06 22:22:31 +01:00
Patrick von Platen 2243a59483 [Euler Discrete] Fix sigma (#6078)
* [Euler Discrete] Fix sigma

* make style
2023-12-06 19:59:38 +01:00
apolinário 466d32c442 [Advanced Diffusion Training] Cache latents to avoid VAE passes for every training step (#6076)
* add cache latents

* style
2023-12-06 14:46:53 +01:00
Dhruv Nair 20ba1fdbbd Disable Tests Fetcher (#6060)
update
2023-12-06 18:10:11 +05:30
Pedro Cuenca ab6672fecd Use CC12M for LCM WDS training example (#5908)
* Fix SD scripts - there are only 2 items per batch

* Adjustments to make the SDXL scripts work with other datasets

* Use public webdataset dataset for examples

* make style

* Minor tweaks to the readmes.

* Stress that the database is illustrative.
2023-12-06 10:35:36 +01:00
Dhruv Nair f90a5139a2 fix 2023-12-06 06:03:58 +00:00
Sayak Paul a2bc2e14b9 [feat] allow SDXL pipeline to run with fused QKV projections (#6030)
* debug

* from step

* print

* turn sigma a list

* make str

* init_noise_sigma

* comment

* remove prints

* feat: introduce fused projections

* change to a better name

* no grad

* device.

* device

* dtype

* okay

* print

* more print

* fix: unbind -> split

* fix: qkv >-> k

* enable disable

* apply attention processor within the method

* attn processors

* _enable_fused_qkv_projections

* remove print

* add fused projection to vae

* add todos.

* add: documentation and cleanups.

* add: test for qkv projection fusion.

* relax assertions.

* relax further

* fix: docs

* fix-copies

* correct error message.

* Empty-Commit

* better conditioning on disable_fused_qkv_projections

* check

* check processor

* bfloat16 computation.

* check latent dtype

* style

* remove copy temporarily

* cast latent to bfloat16

* fix: vae -> self.vae

* remove print.

* add _change_to_group_norm_32

* comment out stuff that didn't work

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* reflect patrick's suggestions.

* fix imports

* fix: disable call.

* fix more

* fix device and dtype

* fix conditions.

* fix more

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-06 07:33:26 +05:30
Arsalan f427345ab1 Device agnostic testing (#5612)
* utils and test modifications to enable device agnostic testing

* device for manual seed in unet1d

* fix generator condition in vae test

* consistency changes to testing

* make style

* add device agnostic testing changes to source and one model test

* make dtype check fns private, log cuda fp16 case

* remove dtype checks from import utils, move to testing_utils

* adding tests for most model classes and one pipeline

* fix vae import
2023-12-05 19:04:13 +05:30
apolinário 6e221334cd [advanced_dreambooth_lora_sdxl_tranining_script] save embeddings locally fix (#6058)
* Update train_dreambooth_lora_sdxl_advanced.py

* remove global function args from dreamboothdataset class

* style

* style

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-05 13:52:34 +01:00
Patrick von Platen 53bc30dd45 [From single file] remove depr warning (#6043) 2023-12-05 18:12:25 +05:30
Radamés Ajna eacf5e34eb Fix demofusion (#6049)
* Update pipeline_demofusion_sdxl.py

* Update README.md
2023-12-05 18:10:46 +05:30
Dhruv Nair 4c05f7856a Ldm unet convert fix (#6038)
* fix

* fix ldm conversion

* fix linting
2023-12-05 18:01:02 +05:30
Dhruv Nair bbd3572044 Pin Ruff Version (#6059)
pinn ruff
2023-12-05 17:51:37 +05:30
Dhruv Nair f948778322 Move kandinsky convert script (#6047)
move kandinsky convert script
2023-12-05 15:12:37 +05:30
Steven Liu 4684ea2fe8 [docs] #Copied from mechanism (#6007)
* copied from section

* feedback
2023-12-04 10:12:52 -08:00
Steven Liu b64f835ea7 [docs] Add Kandinsky 3 (#5988)
* add

* fix api docs

* edits
2023-12-04 10:11:15 -08:00
Linoy Tsaban 880c0fdd36 [advanced dreambooth lora training script][bug_fix] change token_abstraction type to str (#6040)
* improve help tags

* style fix

* changes token_abstraction type to string.
support multiple concepts for pivotal using a comma separated string.

* style fixup

* changed logger to warning (not yet available)

* moved the token_abstraction parsing to be in the same block as where we create the mapping of identifier to token

---------

Co-authored-by: Linoy <linoy@huggingface.co>
2023-12-04 18:38:44 +01:00
RuoyiDu c36f1c3160 [Community Pipeline] DemoFusion: Democratising High-Resolution Image Generation With No $$$ (#6022)
* Add files via upload

* Update README.md

* Update pipeline_demofusion_sdxl.py

* Update pipeline_demofusion_sdxl.py

* Update examples/community/README.md

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-04 19:44:57 +05:30
takuoko 0a08d41961 [Feature] Support IP-Adapter Plus (#5915)
* Support IP-Adapter Plus

* fix format

* restore before black format

* restore before black format

* generic

* Refactor PerceiverAttention

* format

* fix test and refactor PerceiverAttention

* generic encode_image

* keep attention implementation

* merge tests

* encode_image backward compatible

* code quality

* fix controlnet inpaint pipeline

* refactor FFN

* refactor FFN

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
2023-12-04 12:43:34 +01:00
Levi McCallum e185084a5d Add variant argument to dreambooth lora sdxl advanced (#6021)
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-04 12:04:15 +01:00
Dhruv Nair b21729225a Update Tests Fetcher (#5950)
* update setup and deps table

* update

* update

* update

* up

* up

* update

* up

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* update

* quality fix

* fix failure reporting

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-04 12:59:41 +05:30
Parth38 8a812e4e14 Update value_guided_sampling.py (#6027)
* Update value_guided_sampling.py

Changed the scheduler step function as predict_epsilon parameter is not there in latest  DDPM Scheduler

* Update value_guided_sampling.md

Updated a link to a working notebook

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-04 10:36:25 +05:30
gujing bf92e746c0 fix StableDiffusionTensorRT super args error (#6009) 2023-12-04 10:06:23 +05:30
Linoy Tsaban b785a155d6 [advanced dreambooth lora sdxl training script] improve help tags (#6035)
* improve help tags

* style fix

---------

Co-authored-by: Linoy <linoy@huggingface.co>
2023-12-04 09:41:25 +05:30
Sayak Paul d486f0e846 [LoRA serialization] fix: duplicate unet prefix problem. (#5991)
* fix: duplicate unet prefix problem.

* Update src/diffusers/loaders/lora.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-02 21:35:16 +05:30
Sayak Paul 3351270627 [PixArt Tests] remove fast tests from slow suite (#5945)
remove fast tests from slow suite
2023-12-02 20:58:27 +05:30
Junsong Chen 4520e1221a adapt PixArtAlphaPipeline for pixart-lcm model (#5974)
* adapt PixArtAlphaPipeline for pixart-lcm model

* remove original_inference_steps from __call__

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-12-02 13:30:40 +05:30
Long(Tony) Lian 618260409f LLMGroundedDiffusionPipeline: inherit from DiffusionPipeline and fix peft (#6023)
* LLMGroundedDiffusionPipeline: inherit from DiffusionPipeline and fix peft

* Use main in the revision in the examples

* Add "Copied from" statements in comments

* Fix formatting with ruff
2023-12-01 09:58:25 -10:00
Patrick von Platen dadd55fb36 Post Release: v0.24.0 (#5985)
* Post Release: v0.24.0

* post pone deprecation

* post pone deprecation

* Add model_index.json
2023-12-01 18:43:44 +01:00
YiYi Xu 1b6c7ea74e [schedulers] create self.sigmas during __init__ (#6006)
* fix dpm
* all scheulers
2023-12-01 07:15:37 -10:00
YiYi Xu b41f809a4e [Kandinsky 3.0] Follow-up TODOs (#5944)
clean-up kendinsky 3.0
2023-12-01 07:14:22 -10:00
Patrick von Platen 0f55c17e17 fix style 2023-12-01 15:59:34 +00:00
Charchit Sharma 5058d27f12 added attention_head_dim, attention_type, resolution_idx (#6011) 2023-12-01 16:26:58 +01:00
M. Tolga Cangöz 748c1b3ec7 [Docs] Update a link (#6014)
* Update the location of Python's version

* Trim trailing whitespace
2023-12-01 16:26:25 +01:00
M. Tolga Cangöz 523507034f [logging] Fix assertion bug (#6012)
Fix assertion bug
2023-12-01 16:26:04 +01:00
hako-mikan 46c751e970 [Community Pipeline] Regional Prompting Pipeline (#6015)
* Update README.md

* Update README.md

* Add files via upload

* Update README.md

* Update examples/community/README.md

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-12-01 16:22:59 +01:00
Patrick von Platen bc1d28c888 [From Single File] Allow Text Encoder to be passed (#6020)
Allow text encoder to be passed
2023-12-01 16:19:04 +01:00
Sayak Paul af378c1dd1 [Easy] minor edits to setup.py (#5996)
minor edits to setup
2023-12-01 20:38:46 +05:30
Steven Liu 6ba4c5395f [docs] Fix SVD video (#6004)
Update svd.md
2023-12-01 16:07:47 +01:00
Linoy Tsaban c1e4529541 [advanced_dreambooth_lora_sdxl_tranining_script] readme fix (#6019)
readme
2023-12-01 15:14:57 +01:00
Linoy Tsaban d29d97b616 [examples/advanced_diffusion_training] bug fixes and improvements for LoRA Dreambooth SDXL advanced training script (#5935)
* imports and readme bug fixes

* bug fix - ensures text_encoder params are dtype==float32 (when using pivotal tuning) even if the rest of the model is loaded in fp16

* added pivotal tuning to readme

* mapping token identifier to new inserted token in validation prompt (if used)

* correct default value of --train_text_encoder_frac

* change default value of  --adam_weight_decay_text_encoder

* validation prompt generations when using pivotal tuning bug fix

* style fix

* textual inversion embeddings name change

* style fix

* bug fix - stopping text encoder optimization halfway

* readme - will include token abstraction and new inserted tokens when using pivotal tuning
- added type to --num_new_tokens_per_abstraction

* style fix

---------

Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
2023-12-01 14:18:43 +01:00
Jongho Choi 7d4a257c7f Remove a duplicated line? (#6010)
Update __init__.py
2023-12-01 15:49:36 +05:30
Kristian Mischke 141cd52d56 Fix LLMGroundedDiffusionPipeline super class arguments (#5993)
* make `requires_safety_checker` a kwarg instead of a positional argument as it's more future-proof

* apply `make style` formatting edits

* add image_encoder to arguments and pass to super constructor
2023-11-30 10:15:14 -10:00
Steven Liu f72b28c75b [docs] Fix video link (#5986)
Update svd.md
2023-11-29 20:52:25 +01:00
Suraj Patil ada8109d5b Fix SVD doc (#5983)
fix url
2023-11-29 19:55:05 +01:00
Patrick von Platen b34acbdcbc [SDXL Turbo] Add some docs (#5982)
* add diffusers example

* add diffusers example

* Comment about making it faster

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

---------

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2023-11-29 19:52:07 +01:00
Suraj Patil 63f767ef15 Add SVD (#5895)
* begin model

* finish blocks

* add_embedding

* addition_time_embed_dim

* use TimestepEmbedding

* fix temporal res block

* fix time_pos_embed

* fix add_embedding

* add conversion script

* fix model

* up

* add new resnet blocks

* make forward work

* return sample in original shape

* fix temb shape in TemporalResnetBlock

* add spatio temporal transformers

* add vae blocks

* fix blocks

* update

* update

* fix shapes in Alphablender and add time activation in res blcok

* use new blocks

* style

* fix temb shape

* fix SpatioTemporalResBlock

* reuse TemporalBasicTransformerBlock

* fix TemporalBasicTransformerBlock

* use TransformerSpatioTemporalModel

* fix TransformerSpatioTemporalModel

* fix time_context dim

* clean up

* make temb optional

* add blocks

* rename model

* update conversion script

* remove UNetMidBlockSpatioTemporal

* add in init

* remove unused arg

* remove unused arg

* remove more unsed args

* up

* up

* check for None

* update vae

* update up/mid blocks for decoder

* begin pipeline

* adapt scheduler

* add guidance scalings

* fix norm eps in temporal transformers

* add temporal autoencoder

* make pipeline run

* fix frame decodig

* decode in float32

* decode n frames at a time

* pass decoding_t to decode_latents

* fix decode_latents

* vae encode/decode in fp32

* fix dtype in TransformerSpatioTemporalModel

* type image_latents same as image_embeddings

* allow using differnt eps in temporal block for video decoder

* fix default values in vae

* pass num frames in decode

* switch spatial to temporal for mixing in VAE

* fix num frames during split decoding

* cast alpha to sample dtype

* fix attention in MidBlockTemporalDecoder

* fix typo

* fix guidance_scales dtype

* fix missing activation in TemporalDecoder

* skip_post_quant_conv

* add vae conversion

* style

* take guidance scale as input

* up

* allow passing PIL to export_video

* accept fps as arg

* add pipeline and vae in init

* remove hack

* use AutoencoderKLTemporalDecoder

* don't scale image latents

* add unet tests

* clean up unet

* clean TransformerSpatioTemporalModel

* add slow svd test

* clean up

* make temb optional in Decoder mid block

* fix norm eps in TransformerSpatioTemporalModel

* clean up temp decoder

* clean up

* clean up

* use c_noise values for timesteps

* use math for log

* update

* fix copies

* doc

* upcast vae

* update forward pass for gradient checkpointing

* make added_time_ids is tensor

* up

* fix upcasting

* remove post quant conv

* add _resize_with_antialiasing

* fix _compute_padding

* cleanup model

* more cleanup

* more cleanup

* more cleanup

* remove freeu

* remove attn slice

* small clean

* up

* up

* remove extra step kwargs

* remove eta

* remove dropout

* remove callback

* remove merge factor args

* clean

* clean up

* move to dedicated folder

* remove attention_head_dim

* docstr and small fix

* update unet doc strings

* rename decoding_t

* correct linting

* store c_skip and c_out

* cleanup

* clean TemporalResnetBlock

* more cleanup

* clean up vae

* clean up

* begin doc

* more cleanup

* up

* up

* doc

* Improve

* better naming

* better naming

* better naming

* better naming

* better naming

* better naming

* better naming

* better naming

* Apply suggestions from code review

* Default chunk size to None

* add example

* Better

* Apply suggestions from code review

* update doc

* Update src/diffusers/pipelines/stable_diffusion_video/pipeline_stable_diffusion_video.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* style

* Get torch compile working

* up

* rename

* fix doc

* add chunking

* torch compile

* torch compile

* add modelling outputs

* torch compile

* Improve chunking

* Apply suggestions from code review

* Update docs/source/en/using-diffusers/svd.md

* Close diff tag

* remove slicing

* resnet docstr

* add docstr in resnet

* rename

* Apply suggestions from code review

* update tests

* Fix output type latents

* fix more

* fix more

* Update docs/source/en/using-diffusers/svd.md

* fix more

* add pipeline tests

* remove unused arg

* clean  up

* make sure get_scaling receives tensors

* fix euler scheduler

* fix get_scalings

* simply euler for now

* remove old test file

* use randn_tensor to create noise

* fix device for rand tensor

* increase expected_max_difference

* fix test_inference_batch_single_identical

* actually fix test_inference_batch_single_identical

* disable test_save_load_float16

* skip test_float16_inference

* skip test_inference_batch_single_identical

* fix test_xformers_attention_forwardGenerator_pass

* Apply suggestions from code review

* update StableVideoDiffusionPipelineSlowTests

* update image

* add diffusers example

* fix more

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: apolinário <joaopaulo.passos@gmail.com>
2023-11-29 19:13:36 +01:00
PENGUINLIONG d1b2a1a957 Fixed custom module importing on Windows (#5891)
* Fixed custom module importing on Windows

Windows use back slash and `os.path.join()` follows that convention.

* Apply suggestions from code review

Co-authored-by: Lucain <lucainp@gmail.com>

* Update pipeline_utils.py

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Lucain <lucainp@gmail.com>
2023-11-29 16:33:04 +01:00
Kashif Rasul 01782c220e [Wuerstchen] Adapt lora training example scripts to use PEFT (#5959)
* Adapt lora example scripts to use PEFT

* add to_out.0
2023-11-29 16:18:20 +01:00
vahramtadevosyan d63a498c3b [Pipeline] Add TextToVideoZeroSDXLPipeline (#4695)
* integrated sdxl for the text2video-zero pipeline

* make fix-copies

* fixed CI issues

* make fix-copies

* added docs and `copied from` statements

* added fast tests

* made a small change in docs

* quality+style check fix

* updated docs. added controlnet inference with sdxl

* added device compatibility for fast tests

* fixed docstrings

* changing vae upcasting

* remove torch.empty_cache to speed up inference

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* made fast tests to run on dummy models only, fixed copied from statements

* fixed testing utils imports

* Added bullet points for SDXL support

* fixed formatting & quality

* Update tests/pipelines/text_to_video/test_text_to_video_zero_sdxl.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update tests/pipelines/text_to_video/test_text_to_video_zero_sdxl.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* fixed minor error for merging

* fixed updates of sdxl

* made fast tests inherit from `PipelineTesterMixin` and run in 3-4secs on CPU

* make style && make quality

* reimplemented fast tests w/o default attn processor

* make style & make quality

* make fix-copies

* make fix-copies

* fixed docs

* make style & make quality & make fix-copies

* bug fix in cross attention

* make style && make quality

* make fix-copies

* fix gpu issues

* make fix-copies

* updated pipeline signature

---------

Co-authored-by: Vahram <vahram.tadevosyan@lambda-loginnode02.cm.cluster>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2023-11-29 16:10:43 +01:00
Marko Kostiv 6a4aad43dc Controlnet ssd 1b support (#5779)
* Add SSD-1B support for controlnet model

* Add conditioning_channels into ControlNet init from unet

* Fix black formatting

* Isort fixes

* Adds SSD-1B controlnet pipeline test with UNetMidBlock2D as mid block

* Overrides failing ssd-1b tests

* Fixes tests after main branch update

* Fixes code quality checks

---------

Co-authored-by: Marko Kostiv <marko@linearity.io>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-11-29 16:10:01 +01:00
Steven Liu ddd8bd53ed [docs] LCM training (#5796)
* first draft

* feedback
2023-11-29 16:08:05 +01:00
JuanCarlosPi 9f7b2cf2dc Support of ip-adapter to the StableDiffusionControlNetInpaintPipeline (#5887)
* Change pipeline_controlnet_inpaint.py to add ip-adapter support. Changes are similar to those in pipeline_controlnet

* Change tests for the StableDiffusionControlNetInpaintPipeline by adding image_encoder: None

* Update src/diffusers/pipelines/controlnet/pipeline_controlnet_inpaint.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
2023-11-29 16:00:24 +01:00
Sayak Paul 895c4b704b [LoRA refactor] move several state dict conversion utils out of lora.py (#5955)
* move several state dict conversion utils out of lora.py

* check

* check

* check

* check

* check

* check

* check

* revert back

* check

* check

* again check

* maybe fix?

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-29 20:24:21 +05:30
Linh Nguyen 636feba552 Rename output_dir argument (#5916)
Fix typo in output_dir argument: "text-inversion-model" → "dreambooth-model"
2023-11-29 15:47:16 +01:00
Andrés Romero 79dc7df03e [bug fix] Inpainting for MultiAdapter (#5922)
* bug in MultiAdapter for Inpainting

* adapter_input is a list for MultiAdapter

---------

Co-authored-by: andres <andres@hax.ai>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-11-29 15:46:26 +01:00
Charchit Sharma 6031ecbd23 added doc for Kandinsky3.0 (#5937)
* added en doc for Kandinsky3.0

* required changes

* Update docs/source/en/api/pipelines/kandinsky3.md

* Update docs/source/en/api/pipelines/kandinsky3.md

* Update docs/source/en/api/pipelines/kandinsky3.md

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-29 15:32:00 +01:00
Sayak Paul fdd003d8e2 [Tests] Refactor test_examples.py for better readability (#5946)
* control and custom diffusion

* dreambooth

* instructpix2pix and dreambooth ckpting

* t2i adapters.

* text to image ft

* textual inversion

* unconditional

* workflows

* import fix

* fix import
2023-11-29 18:43:59 +05:30
Steven Liu 172acc98b9 [docs] Update pipeline list (#5952)
add to list
2023-11-29 14:08:39 +01:00
estelleafl 5ae3c3a56b [ldm3d] Ldm3d upscaler to community pipeline (#5870)
---------
Co-authored-by: Aflalo <estellea@isl-gpu27.rr.intel.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
2023-11-28 09:00:39 -10:00
Soumik Rakshit 21bc59ab24 fix: minor typo in docstring (#5961) 2023-11-28 18:18:34 +05:30
Steven Liu 50a749e909 [docs] Fix space (#5898)
* fix

* minor edits
2023-11-27 11:50:59 -08:00
YiYi Xu d9075be494 [load_textual_inversion]: allow multiple tokens (#5837)
Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-11-27 06:52:36 -10:00
Patrick von Platen b135b6e905 [From_pretrained] Fix warning (#5948) 2023-11-27 14:35:19 +01:00
T. Xu 14a0d21d2e [Community Pipeline] Diffusion Posterior Sampling for General Noisy Inverse Problems (#5939)
* [community pipeline] dps impl

* add type checking

* pass ruff check

* ruff formatter
2023-11-27 14:29:42 +01:00
Patrick von Platen ebf581e85f [Tests] Make sure that we don't run tests multiple times (#5949)
* [Tests] Make sure that we don't run tests mulitple times

* [Tests] Make sure that we don't run tests mulitple times

* [Tests] Make sure that we don't run tests mulitple times
2023-11-27 14:18:56 +01:00
Patrick von Platen e550163b9f [Vae] Make sure all vae's work with latent diffusion models (#5880)
* add comments to explain the code better

* add comments to explain the code better

* add comments to explain the code better

* add comments to explain the code better

* add comments to explain the code better

* fix more

* fix more

* fix more

* fix more

* fix more

* fix more
2023-11-27 14:17:47 +01:00
Viktor Grygorchuk 20f0cbc88f fix: error on device for lpw_stable_diffusion_xl pipeline if pipe.enable_sequential_cpu_offload() enabled (#5885)
fix: set device for pipe.enable_sequential_cpu_offload()
2023-11-27 13:47:47 +01:00
Chi d72a24b790 Replace multiple variables with one variable. (#5715)
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.

* Update src/diffusers/models/unet_2d_blocks.py

This changes suggest by maintener.

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update src/diffusers/models/unet_2d_blocks.py

Add suggested text

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update unet_2d_blocks.py

I changed the Parameter to Args text.

* Update unet_2d_blocks.py

proper indentation set in this file.

* Update unet_2d_blocks.py

a little bit of change in the act_fun argument line.

* I run the black command to reformat style in the code

* Update unet_2d_blocks.py

similar doc-string add to have in the original diffusion repository.

* I enhanced the code by replacing multiple redundant variables with a single variable, as they all served the same purpose. Additionally, I utilized the get_activation function for improved flexibility in choosing activation functions.

* Using as black package to reformated my file

* reverte some changes

* Remove conv_out_padding variables and using as conv_in_padding

* conv_out_padding create and add them into the code.

* run black command to solving styling problem

* add little bit space between comment and import statement

* I am utilizing the ruff library to address the style issues in my Makefile.

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-27 13:34:52 +01:00
ginjia d3cda804e7 add LoRA weights load and fuse support for IPEX pipeline (#5920)
add IPEX pipeline LoRA weights loading support
2023-11-27 13:32:43 +01:00
dg845 07eac4d65a Fix LCM Stable Diffusion distillation bug related to parsing unet_time_cond_proj_dim (#5893)
* Fix bug related to parsing unet_time_cond_proj_dim.

* Fix analogous bug in the SD-XL LCM distillation script.
2023-11-27 13:00:40 +01:00
Iván de Prado c079cae3d4 Avoid computing min() that is expensive when do_normalize is False in the image processor (#5896)
Avoid computing min() that is expensive when do_normalize is False

Avoid extra computing when do_normalize is False
2023-11-27 12:46:26 +01:00
Wang, Yi c7bfb8b22a set the model to train state before accelerator prepare (#5099)
Signed-off-by: Wang, Yi A <yi.a.wang@intel.com>
2023-11-27 12:43:49 +01:00
dg845 67d070749a Add Custom Timesteps Support to LCMScheduler and Supported Pipelines (#5874)
* Add custom timesteps support to LCMScheduler.

* Add custom timesteps support to StableDiffusionPipeline.

* Add custom timesteps support to StableDiffusionXLPipeline.

* Add custom timesteps support to remaining Stable Diffusion pipelines which support LCMScheduler (img2img, inpaint).

* Add custom timesteps support to remaining Stable Diffusion XL pipelines which support LCMScheduler (img2img, inpaint).

* Add custom timesteps support to StableDiffusionControlNetPipeline.

* Add custom timesteps support to T21 Stable Diffusion (XL) Adapters.

* Clean up Stable Diffusion inpaint tests.

* Manually add support for custom timesteps to AltDiffusion pipelines since make fix-copies doesn't appear to work correctly (it deletes the whole pipeline).

* make style

* Refactor pipeline timestep handling into the retrieve_timesteps function.
2023-11-27 12:39:14 +01:00
Aryan V S 9c357bda3f Deprecate KarrasVeScheduler and ScoreSdeVpScheduler (#5269)
* deprecated: KarrasVeScheduler, ScoreSdeVpScheduler

* delete tests relevant to deprecated schedulers

* chore: run make style

* fix: import error caused due to incorrect _import_structure after deprecation

* fix: ScoreSdeVpScheduler was not importable from diffusers

* remove import added by assumption

* Update src/diffusers/schedulers/__init__.py as suggested by @patrickvonplaten

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* make it a part deprecated

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Fix

* fix

* fix doc

* fix doc....again.......

* remove karras_ve test folder

Co-Authored-By: YiYi Xu <yixu310@gmail.com>

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-11-27 12:33:02 +01:00
Sayak Paul 3f7c3511dc [Core] add support for gradient checkpointing in transformer_2d (#5943)
add support for gradient checkpointing in transformer_2d
2023-11-27 16:21:12 +05:30
Junsong Chen 7d6f30e89b [Fix: pixart-alpha] random 512px resolution bug (#5842)
* [Fix: pixart-alpha]
add ASPECT_RATIO_512_BIN in use_resolution_binning for random 512px image generation.

* add slow test file for 512px generation without resolution binning

* fix: slow tests for resolution binning.

---------

Co-authored-by: jschen <chenjunsong4@h-partners.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2023-11-27 13:35:35 +05:30
Patrick von Platen 6d2e19f746 [Examples] Allow downloading variant model files (#5531)
* add variant

* add variant

* Apply suggestions from code review

* reformat

* fix: textual_inversion.py

* fix: variant in model_info

---------

Co-authored-by: sayakpaul <spsayakpaul@gmail.com>
2023-11-27 10:43:20 +05:30
Patrick von Platen 2a7f43a73b correct num inference steps 2023-11-24 17:09:26 +00:00
Patrick von Platen b978334d71 [@cene555][Kandinsky 3.0] Add Kandinsky 3.0 (#5913)
* finalize

* finalize

* finalize

* add slow test

* add slow test

* add slow test

* Fix more

* add slow test

* fix more

* fix more

* fix more

* fix more

* fix more

* fix more

* fix more

* fix more

* fix more

* Better

* Fix more

* Fix more

* add slow test

* Add auto pipelines

* add slow test

* Add all

* add slow test

* add slow test

* add slow test

* add slow test

* add slow test

* Apply suggestions from code review

* add slow test

* add slow test
2023-11-24 17:46:00 +01:00
Sayak Paul e5f232f76b [Docs] add: 8bit inference with pixart alpha (#5814)
* add: 8bit inference with pixart alpha

* Apply suggestions from code review

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* add: note on 4bit.

* Apply suggestions from code review

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* address comment

---------

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-24 20:36:33 +05:30
Linoy Tsaban 3003ff4947 [bug fix] fix small bug in readme template of sdxl lora training script (#5914)
readme improvement and metadata fix
2023-11-23 19:08:49 +01:00
Linoy Tsaban 5ffa603244 [bug fix] fix small bug in readme template of sdxl lora training script (#5906)
* readme bug fix

* style fix

---------

Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
2023-11-23 12:11:50 +01:00
Linoy Tsaban 0eeee618cf Adds an advanced version of the SD-XL DreamBooth LoRA training script supporting pivotal tuning (#5883)
* sdxl dreambooth lora training script with pivotal tuning

* bug fix - args missing from parse_args

* code quality fixes

* comment unnecessary code from TokenEmbedding handler class

* fixup

---------

Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
2023-11-22 16:27:56 +01:00
Andrés Romero 93f1a14cab ControlNet+Adapter pipeline, and ControlNet+Adapter+Inpaint pipeline (#5869)
* ControlNet+Adapter pipeline, and +Inpaint pipeline


---------

Co-authored-by: andres <andres@hax.ai>
2023-11-21 08:59:29 -10:00
Patrick von Platen 13d73d9303 [Lora] Seperate logic (#5809)
* [Lora] Seperate logic

* [Lora] Seperate logic

* [Lora] Seperate logic

* add comments to explain the code better

* add comments to explain the code better
2023-11-21 18:58:37 +01:00
YiYi Xu ba352aea29 [feat] IP Adapters (author @okotaku ) (#5713)
* add ip-adapter


---------

Co-authored-by: okotaku <to78314910@gmail.com>
Co-authored-by: sayakpaul <spsayakpaul@gmail.com>
Co-authored-by: yiyixuxu <yixu310@gmail,com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
2023-11-21 07:34:30 -10:00
Linoy Tsaban 6fac1369d0 Add features to the Dreambooth LoRA SDXL training script (#5508)
* Additions:
- support for different lr for text encoder
- support for Prodigy optimizer
- support for min snr gamma
- support for custom captions and dataset loading from the hub

* adjusted --caption_column behaviour (to -not- use the second column of the dataset by default if --caption_column is not provided)

* fixed --output_dir / --model_dir_name confusion

* added --repeats, --adam_weight_decay_text_encoder
+ some fixes

* Update examples/dreambooth/train_dreambooth_lora_sdxl.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update examples/dreambooth/train_dreambooth_lora_sdxl.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update examples/dreambooth/train_dreambooth_lora_sdxl.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* - import compute_snr from diffusers/training_utils.py
- cluster adamw together
- when using 'prodigy', if --train_text_encoder == True and --text_encoder_lr != --learning rate, changes the lr of the text encoders optimization params to be --learning_rate (otherwise errors)

* shape fixes when custom captions are used

* formatting and a little cleanup

* code styling

* --repeats default value fixed, changed to 1

* bug fix - removed redundant lines of embedding concatenation when using prior_preservation (that duplicated class_prompt embeddings)

* changed dataset loading logic according to the following usecases (to avoid unnecessary dependency on datasets)-
1. user provides --dataset_name
2. user provides local dir --instance_data_dir that contains a metadata .jsonl file
3. user provides local dir --instance_data_dir that contains only images
in cases [1,2] we import datasets and use load_dataset method, in case [3] we process the data same as in the original script setting

* styling fix

* arg name fix

* adjusted the --repeats logic

* -removed redundant arg and 'if' when loading local folder with prompts
-updated readme template
-some default val fixes
-custom caption tests

* image path fix for readme

* code style

* bug fix

* --caption_column arg

* readme fix

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
2023-11-21 17:38:43 +01:00
Steven Liu 1093f9d615 [docs] MusicLDM (#5854)
* fix

* feedback
2023-11-21 15:27:41 +01:00
Aryan V S 81780882b8 Addition of new callbacks to controlnets (#5812)
* add new callbacks to src/diffusers/pipelines/controlnet/pipeline_controlnet.py

* update callbacks

* fix repeated kwarg

* update

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2023-11-21 15:22:20 +01:00
Dhruv Nair ebc7bedeb7 Add tests fetcher (#5848)
* add tests fetcher to utils

* add test fetcher

* update

* update

* remove unused dependency version check script

* update

* fix mistake

* update

* update

* update

* update

* update

* update

* update

* remove concurrency params

* update

* update

* update

* update

* update

* update

* move test fetcher to dedicated workflow
2023-11-21 18:01:44 +05:30
306 changed files with 37886 additions and 4136 deletions
+52
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@@ -0,0 +1,52 @@
name: Benchmarking tests
on:
schedule:
- cron: "30 1 1,15 * *" # every 2 weeks on the 1st and the 15th of every month at 1:30 AM
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
jobs:
torch_pipelines_cuda_benchmark_tests:
name: Torch Core Pipelines CUDA Benchmarking Tests
strategy:
fail-fast: false
max-parallel: 1
runs-on: [single-gpu, nvidia-gpu, a10, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install pandas
- name: Environment
run: |
python utils/print_env.py
- name: Diffusers Benchmarking
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.DIFFUSERS_BOT_TOKEN }}
BASE_PATH: benchmark_outputs
run: |
export TOTAL_GPU_MEMORY=$(python -c "import torch; print(torch.cuda.get_device_properties(0).total_memory / (1024**3))")
cd benchmarks && mkdir ${BASE_PATH} && python run_all.py && python push_results.py
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: benchmark_test_reports
path: benchmarks/benchmark_outputs
+170
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@@ -0,0 +1,170 @@
name: Fast tests for PRs - Test Fetcher
on: workflow_dispatch
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
setup_pr_tests:
name: Setup PR Tests
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
outputs:
matrix: ${{ steps.set_matrix.outputs.matrix }}
test_map: ${{ steps.set_matrix.outputs.test_map }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 0
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python utils/print_env.py
echo $(git --version)
- name: Fetch Tests
run: |
python utils/tests_fetcher.py | tee test_preparation.txt
- name: Report fetched tests
uses: actions/upload-artifact@v3
with:
name: test_fetched
path: test_preparation.txt
- id: set_matrix
name: Create Test Matrix
# The `keys` is used as GitHub actions matrix for jobs, i.e. `models`, `pipelines`, etc.
# The `test_map` is used to get the actual identified test files under each key.
# If no test to run (so no `test_map.json` file), create a dummy map (empty matrix will fail)
run: |
if [ -f test_map.json ]; then
keys=$(python3 -c 'import json; fp = open("test_map.json"); test_map = json.load(fp); fp.close(); d = list(test_map.keys()); print(json.dumps(d))')
test_map=$(python3 -c 'import json; fp = open("test_map.json"); test_map = json.load(fp); fp.close(); print(json.dumps(test_map))')
else
keys=$(python3 -c 'keys = ["dummy"]; print(keys)')
test_map=$(python3 -c 'test_map = {"dummy": []}; print(test_map)')
fi
echo $keys
echo $test_map
echo "matrix=$keys" >> $GITHUB_OUTPUT
echo "test_map=$test_map" >> $GITHUB_OUTPUT
run_pr_tests:
name: Run PR Tests
needs: setup_pr_tests
if: contains(fromJson(needs.setup_pr_tests.outputs.matrix), 'dummy') != true
strategy:
fail-fast: false
max-parallel: 2
matrix:
modules: ${{ fromJson(needs.setup_pr_tests.outputs.matrix) }}
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install accelerate
- name: Environment
run: |
python utils/print_env.py
- name: Run all selected tests on CPU
run: |
python -m pytest -n 2 --dist=loadfile -v --make-reports=${{ matrix.modules }}_tests_cpu ${{ fromJson(needs.setup_pr_tests.outputs.test_map)[matrix.modules] }}
- name: Failure short reports
if: ${{ failure() }}
continue-on-error: true
run: |
cat reports/${{ matrix.modules }}_tests_cpu_stats.txt
cat reports/${{ matrix.modules }}_tests_cpu_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v3
with:
name: ${{ matrix.modules }}_test_reports
path: reports
run_staging_tests:
strategy:
fail-fast: false
matrix:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Run Hub tests for models, schedulers, and pipelines on a staging env
if: ${{ matrix.config.framework == 'hub_tests_pytorch' }}
run: |
HUGGINGFACE_CO_STAGING=true python -m pytest \
-m "is_staging_test" \
--make-reports=tests_${{ matrix.config.report }} \
tests
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports
+2 -1
View File
@@ -113,9 +113,10 @@ jobs:
- name: Run example PyTorch CPU tests
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples/test_examples.py
examples
- name: Failure short reports
if: ${{ failure() }}
+5 -1
View File
@@ -5,6 +5,10 @@ on:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
@@ -96,7 +100,7 @@ jobs:
run: |
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples/test_examples.py
examples
- name: Failure short reports
if: ${{ failure() }}
+4
View File
@@ -13,6 +13,10 @@ env:
PYTEST_TIMEOUT: 600
RUN_SLOW: no
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
run_fast_tests_apple_m1:
name: Fast PyTorch MPS tests on MacOS
+1 -1
View File
@@ -355,7 +355,7 @@ You will need basic `git` proficiency to be able to contribute to
manual. Type `git --help` in a shell and enjoy. If you prefer books, [Pro
Git](https://git-scm.com/book/en/v2) is a very good reference.
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/main/setup.py#L244)):
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/main/setup.py#L265)):
1. Fork the [repository](https://github.com/huggingface/diffusers) by
clicking on the 'Fork' button on the repository's page. This creates a copy of the code
+2 -2
View File
@@ -3,7 +3,7 @@
# make sure to test the local checkout in scripts and not the pre-installed one (don't use quotes!)
export PYTHONPATH = src
check_dirs := examples scripts src tests utils
check_dirs := examples scripts src tests utils benchmarks
modified_only_fixup:
$(eval modified_py_files := $(shell python utils/get_modified_files.py $(check_dirs)))
@@ -41,7 +41,7 @@ repo-consistency:
quality:
ruff check $(check_dirs) setup.py
ruff format --check $(check_dirs) setup.py
ruff format --check $(check_dirs) setup.py
python utils/check_doc_toc.py
# Format source code automatically and check is there are any problems left that need manual fixing
+1 -1
View File
@@ -82,7 +82,7 @@ Models are designed as configurable toolboxes that are natural extensions of [Py
The following design principles are followed:
- Models correspond to **a type of model architecture**. *E.g.* the [`UNet2DConditionModel`] class is used for all UNet variations that expect 2D image inputs and are conditioned on some context.
- All models can be found in [`src/diffusers/models`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and every model architecture shall be defined in its file, e.g. [`unet_2d_condition.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_condition.py), [`transformer_2d.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/transformer_2d.py), etc...
- Models **do not** follow the single-file policy and should make use of smaller model building blocks, such as [`attention.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention.py), [`resnet.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/resnet.py), [`embeddings.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py), etc... **Note**: This is in stark contrast to Transformers' modelling files and shows that models do not really follow the single-file policy.
- Models **do not** follow the single-file policy and should make use of smaller model building blocks, such as [`attention.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention.py), [`resnet.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/resnet.py), [`embeddings.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py), etc... **Note**: This is in stark contrast to Transformers' modeling files and shows that models do not really follow the single-file policy.
- Models intend to expose complexity, just like PyTorch's `Module` class, and give clear error messages.
- Models all inherit from `ModelMixin` and `ConfigMixin`.
- Models can be optimized for performance when it doesnt demand major code changes, keep backward compatibility, and give significant memory or compute gain.
+2 -2
View File
@@ -77,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 15000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 16000+ checkpoints):
```python
from diffusers import DiffusionPipeline
@@ -219,7 +219,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +6000 other amazing GitHub repositories 💪
- +7000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
+297
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@@ -0,0 +1,297 @@
import os
import sys
import torch
from diffusers import (
AutoPipelineForImage2Image,
AutoPipelineForInpainting,
AutoPipelineForText2Image,
ControlNetModel,
LCMScheduler,
StableDiffusionAdapterPipeline,
StableDiffusionControlNetPipeline,
StableDiffusionXLAdapterPipeline,
StableDiffusionXLControlNetPipeline,
T2IAdapter,
WuerstchenCombinedPipeline,
)
from diffusers.utils import load_image
sys.path.append(".")
from utils import ( # noqa: E402
BASE_PATH,
PROMPT,
BenchmarkInfo,
benchmark_fn,
bytes_to_giga_bytes,
flush,
generate_csv_dict,
write_to_csv,
)
RESOLUTION_MAPPING = {
"runwayml/stable-diffusion-v1-5": (512, 512),
"lllyasviel/sd-controlnet-canny": (512, 512),
"diffusers/controlnet-canny-sdxl-1.0": (1024, 1024),
"TencentARC/t2iadapter_canny_sd14v1": (512, 512),
"TencentARC/t2i-adapter-canny-sdxl-1.0": (1024, 1024),
"stabilityai/stable-diffusion-2-1": (768, 768),
"stabilityai/stable-diffusion-xl-base-1.0": (1024, 1024),
"stabilityai/stable-diffusion-xl-refiner-1.0": (1024, 1024),
"stabilityai/sdxl-turbo": (512, 512),
}
class BaseBenchmak:
pipeline_class = None
def __init__(self, args):
super().__init__()
def run_inference(self, args):
raise NotImplementedError
def benchmark(self, args):
raise NotImplementedError
def get_result_filepath(self, args):
pipeline_class_name = str(self.pipe.__class__.__name__)
name = (
args.ckpt.replace("/", "_")
+ "_"
+ pipeline_class_name
+ f"-bs@{args.batch_size}-steps@{args.num_inference_steps}-mco@{args.model_cpu_offload}-compile@{args.run_compile}.csv"
)
filepath = os.path.join(BASE_PATH, name)
return filepath
class TextToImageBenchmark(BaseBenchmak):
pipeline_class = AutoPipelineForText2Image
def __init__(self, args):
pipe = self.pipeline_class.from_pretrained(args.ckpt, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
if args.run_compile:
if not isinstance(pipe, WuerstchenCombinedPipeline):
pipe.unet.to(memory_format=torch.channels_last)
print("Run torch compile")
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
if hasattr(pipe, "movq") and getattr(pipe, "movq", None) is not None:
pipe.movq.to(memory_format=torch.channels_last)
pipe.movq = torch.compile(pipe.movq, mode="reduce-overhead", fullgraph=True)
else:
print("Run torch compile")
pipe.decoder = torch.compile(pipe.decoder, mode="reduce-overhead", fullgraph=True)
pipe.vqgan = torch.compile(pipe.vqgan, mode="reduce-overhead", fullgraph=True)
pipe.set_progress_bar_config(disable=True)
self.pipe = pipe
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
)
def benchmark(self, args):
flush()
print(f"[INFO] {self.pipe.__class__.__name__}: Running benchmark with: {vars(args)}\n")
time = benchmark_fn(self.run_inference, self.pipe, args) # in seconds.
memory = bytes_to_giga_bytes(torch.cuda.max_memory_allocated()) # in GBs.
benchmark_info = BenchmarkInfo(time=time, memory=memory)
pipeline_class_name = str(self.pipe.__class__.__name__)
flush()
csv_dict = generate_csv_dict(
pipeline_cls=pipeline_class_name, ckpt=args.ckpt, args=args, benchmark_info=benchmark_info
)
filepath = self.get_result_filepath(args)
write_to_csv(filepath, csv_dict)
print(f"Logs written to: {filepath}")
flush()
class TurboTextToImageBenchmark(TextToImageBenchmark):
def __init__(self, args):
super().__init__(args)
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
guidance_scale=0.0,
)
class LCMLoRATextToImageBenchmark(TextToImageBenchmark):
lora_id = "latent-consistency/lcm-lora-sdxl"
def __init__(self, args):
super().__init__(args)
self.pipe.load_lora_weights(self.lora_id)
self.pipe.fuse_lora()
self.pipe.scheduler = LCMScheduler.from_config(self.pipe.scheduler.config)
def get_result_filepath(self, args):
pipeline_class_name = str(self.pipe.__class__.__name__)
name = (
self.lora_id.replace("/", "_")
+ "_"
+ pipeline_class_name
+ f"-bs@{args.batch_size}-steps@{args.num_inference_steps}-mco@{args.model_cpu_offload}-compile@{args.run_compile}.csv"
)
filepath = os.path.join(BASE_PATH, name)
return filepath
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
guidance_scale=1.0,
)
class ImageToImageBenchmark(TextToImageBenchmark):
pipeline_class = AutoPipelineForImage2Image
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/1665_Girl_with_a_Pearl_Earring.jpg"
image = load_image(url).convert("RGB")
def __init__(self, args):
super().__init__(args)
self.image = self.image.resize(RESOLUTION_MAPPING[args.ckpt])
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
image=self.image,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
)
class TurboImageToImageBenchmark(ImageToImageBenchmark):
def __init__(self, args):
super().__init__(args)
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
image=self.image,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
guidance_scale=0.0,
strength=0.5,
)
class InpaintingBenchmark(ImageToImageBenchmark):
pipeline_class = AutoPipelineForInpainting
mask_url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/overture-creations-5sI6fQgYIuo_mask.png"
mask = load_image(mask_url).convert("RGB")
def __init__(self, args):
super().__init__(args)
self.image = self.image.resize(RESOLUTION_MAPPING[args.ckpt])
self.mask = self.mask.resize(RESOLUTION_MAPPING[args.ckpt])
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
image=self.image,
mask_image=self.mask,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
)
class ControlNetBenchmark(TextToImageBenchmark):
pipeline_class = StableDiffusionControlNetPipeline
aux_network_class = ControlNetModel
root_ckpt = "runwayml/stable-diffusion-v1-5"
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/canny_image_condition.png"
image = load_image(url).convert("RGB")
def __init__(self, args):
aux_network = self.aux_network_class.from_pretrained(args.ckpt, torch_dtype=torch.float16)
pipe = self.pipeline_class.from_pretrained(self.root_ckpt, controlnet=aux_network, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
pipe.set_progress_bar_config(disable=True)
self.pipe = pipe
if args.run_compile:
pipe.unet.to(memory_format=torch.channels_last)
pipe.controlnet.to(memory_format=torch.channels_last)
print("Run torch compile")
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipe.controlnet = torch.compile(pipe.controlnet, mode="reduce-overhead", fullgraph=True)
self.image = self.image.resize(RESOLUTION_MAPPING[args.ckpt])
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
image=self.image,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
)
class ControlNetSDXLBenchmark(ControlNetBenchmark):
pipeline_class = StableDiffusionXLControlNetPipeline
root_ckpt = "stabilityai/stable-diffusion-xl-base-1.0"
def __init__(self, args):
super().__init__(args)
class T2IAdapterBenchmark(ControlNetBenchmark):
pipeline_class = StableDiffusionAdapterPipeline
aux_network_class = T2IAdapter
root_ckpt = "CompVis/stable-diffusion-v1-4"
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/canny_for_adapter.png"
image = load_image(url).convert("L")
def __init__(self, args):
aux_network = self.aux_network_class.from_pretrained(args.ckpt, torch_dtype=torch.float16)
pipe = self.pipeline_class.from_pretrained(self.root_ckpt, adapter=aux_network, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
pipe.set_progress_bar_config(disable=True)
self.pipe = pipe
if args.run_compile:
pipe.unet.to(memory_format=torch.channels_last)
pipe.adapter.to(memory_format=torch.channels_last)
print("Run torch compile")
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipe.adapter = torch.compile(pipe.adapter, mode="reduce-overhead", fullgraph=True)
self.image = self.image.resize(RESOLUTION_MAPPING[args.ckpt])
class T2IAdapterSDXLBenchmark(T2IAdapterBenchmark):
pipeline_class = StableDiffusionXLAdapterPipeline
root_ckpt = "stabilityai/stable-diffusion-xl-base-1.0"
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/canny_for_adapter_sdxl.png"
image = load_image(url)
def __init__(self, args):
super().__init__(args)
+26
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@@ -0,0 +1,26 @@
import argparse
import sys
sys.path.append(".")
from base_classes import ControlNetBenchmark, ControlNetSDXLBenchmark # noqa: E402
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="lllyasviel/sd-controlnet-canny",
choices=["lllyasviel/sd-controlnet-canny", "diffusers/controlnet-canny-sdxl-1.0"],
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_pipe = (
ControlNetBenchmark(args) if args.ckpt == "lllyasviel/sd-controlnet-canny" else ControlNetSDXLBenchmark(args)
)
benchmark_pipe.benchmark(args)
+29
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@@ -0,0 +1,29 @@
import argparse
import sys
sys.path.append(".")
from base_classes import ImageToImageBenchmark, TurboImageToImageBenchmark # noqa: E402
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="runwayml/stable-diffusion-v1-5",
choices=[
"runwayml/stable-diffusion-v1-5",
"stabilityai/stable-diffusion-2-1",
"stabilityai/stable-diffusion-xl-refiner-1.0",
"stabilityai/sdxl-turbo",
],
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_pipe = ImageToImageBenchmark(args) if "turbo" not in args.ckpt else TurboImageToImageBenchmark(args)
benchmark_pipe.benchmark(args)
+28
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@@ -0,0 +1,28 @@
import argparse
import sys
sys.path.append(".")
from base_classes import InpaintingBenchmark # noqa: E402
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="runwayml/stable-diffusion-v1-5",
choices=[
"runwayml/stable-diffusion-v1-5",
"stabilityai/stable-diffusion-2-1",
"stabilityai/stable-diffusion-xl-base-1.0",
],
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_pipe = InpaintingBenchmark(args)
benchmark_pipe.benchmark(args)
+28
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@@ -0,0 +1,28 @@
import argparse
import sys
sys.path.append(".")
from base_classes import T2IAdapterBenchmark, T2IAdapterSDXLBenchmark # noqa: E402
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="TencentARC/t2iadapter_canny_sd14v1",
choices=["TencentARC/t2iadapter_canny_sd14v1", "TencentARC/t2i-adapter-canny-sdxl-1.0"],
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_pipe = (
T2IAdapterBenchmark(args)
if args.ckpt == "TencentARC/t2iadapter_canny_sd14v1"
else T2IAdapterSDXLBenchmark(args)
)
benchmark_pipe.benchmark(args)
+23
View File
@@ -0,0 +1,23 @@
import argparse
import sys
sys.path.append(".")
from base_classes import LCMLoRATextToImageBenchmark # noqa: E402
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="stabilityai/stable-diffusion-xl-base-1.0",
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=4)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_pipe = LCMLoRATextToImageBenchmark(args)
benchmark_pipe.benchmark(args)
+40
View File
@@ -0,0 +1,40 @@
import argparse
import sys
sys.path.append(".")
from base_classes import TextToImageBenchmark, TurboTextToImageBenchmark # noqa: E402
ALL_T2I_CKPTS = [
"runwayml/stable-diffusion-v1-5",
"segmind/SSD-1B",
"stabilityai/stable-diffusion-xl-base-1.0",
"kandinsky-community/kandinsky-2-2-decoder",
"warp-ai/wuerstchen",
"stabilityai/sdxl-turbo",
]
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="runwayml/stable-diffusion-v1-5",
choices=ALL_T2I_CKPTS,
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
benchmark_cls = None
if "turbo" in args.ckpt:
benchmark_cls = TurboTextToImageBenchmark
else:
benchmark_cls = TextToImageBenchmark
benchmark_pipe = benchmark_cls(args)
benchmark_pipe.benchmark(args)
+72
View File
@@ -0,0 +1,72 @@
import glob
import sys
import pandas as pd
from huggingface_hub import hf_hub_download, upload_file
from huggingface_hub.utils._errors import EntryNotFoundError
sys.path.append(".")
from utils import BASE_PATH, FINAL_CSV_FILE, GITHUB_SHA, REPO_ID, collate_csv # noqa: E402
def has_previous_benchmark() -> str:
csv_path = None
try:
csv_path = hf_hub_download(repo_id=REPO_ID, repo_type="dataset", filename=FINAL_CSV_FILE)
except EntryNotFoundError:
csv_path = None
return csv_path
def filter_float(value):
if isinstance(value, str):
return float(value.split()[0])
return value
def push_to_hf_dataset():
all_csvs = sorted(glob.glob(f"{BASE_PATH}/*.csv"))
collate_csv(all_csvs, FINAL_CSV_FILE)
# If there's an existing benchmark file, we should report the changes.
csv_path = has_previous_benchmark()
if csv_path is not None:
current_results = pd.read_csv(FINAL_CSV_FILE)
previous_results = pd.read_csv(csv_path)
numeric_columns = current_results.select_dtypes(include=["float64", "int64"]).columns
numeric_columns = [
c for c in numeric_columns if c not in ["batch_size", "num_inference_steps", "actual_gpu_memory (gbs)"]
]
for column in numeric_columns:
previous_results[column] = previous_results[column].map(lambda x: filter_float(x))
# Calculate the percentage change
current_results[column] = current_results[column].astype(float)
previous_results[column] = previous_results[column].astype(float)
percent_change = ((current_results[column] - previous_results[column]) / previous_results[column]) * 100
# Format the values with '+' or '-' sign and append to original values
current_results[column] = current_results[column].map(str) + percent_change.map(
lambda x: f" ({'+' if x > 0 else ''}{x:.2f}%)"
)
# There might be newly added rows. So, filter out the NaNs.
current_results[column] = current_results[column].map(lambda x: x.replace(" (nan%)", ""))
# Overwrite the current result file.
current_results.to_csv(FINAL_CSV_FILE, index=False)
commit_message = f"upload from sha: {GITHUB_SHA}" if GITHUB_SHA is not None else "upload benchmark results"
upload_file(
repo_id=REPO_ID,
path_in_repo=FINAL_CSV_FILE,
path_or_fileobj=FINAL_CSV_FILE,
repo_type="dataset",
commit_message=commit_message,
)
if __name__ == "__main__":
push_to_hf_dataset()
+97
View File
@@ -0,0 +1,97 @@
import glob
import subprocess
import sys
from typing import List
sys.path.append(".")
from benchmark_text_to_image import ALL_T2I_CKPTS # noqa: E402
PATTERN = "benchmark_*.py"
class SubprocessCallException(Exception):
pass
# Taken from `test_examples_utils.py`
def run_command(command: List[str], return_stdout=False):
"""
Runs `command` with `subprocess.check_output` and will potentially return the `stdout`. Will also properly capture
if an error occurred while running `command`
"""
try:
output = subprocess.check_output(command, stderr=subprocess.STDOUT)
if return_stdout:
if hasattr(output, "decode"):
output = output.decode("utf-8")
return output
except subprocess.CalledProcessError as e:
raise SubprocessCallException(
f"Command `{' '.join(command)}` failed with the following error:\n\n{e.output.decode()}"
) from e
def main():
python_files = glob.glob(PATTERN)
for file in python_files:
print(f"****** Running file: {file} ******")
# Run with canonical settings.
if file != "benchmark_text_to_image.py":
command = f"python {file}"
run_command(command.split())
command += " --run_compile"
run_command(command.split())
# Run variants.
for file in python_files:
if file == "benchmark_text_to_image.py":
for ckpt in ALL_T2I_CKPTS:
command = f"python {file} --ckpt {ckpt}"
if "turbo" in ckpt:
command += " --num_inference_steps 1"
run_command(command.split())
command += " --run_compile"
run_command(command.split())
elif file == "benchmark_sd_img.py":
for ckpt in ["stabilityai/stable-diffusion-xl-refiner-1.0", "stabilityai/sdxl-turbo"]:
command = f"python {file} --ckpt {ckpt}"
if ckpt == "stabilityai/sdxl-turbo":
command += " --num_inference_steps 2"
run_command(command.split())
command += " --run_compile"
run_command(command.split())
elif file == "benchmark_sd_inpainting.py":
sdxl_ckpt = "stabilityai/stable-diffusion-xl-base-1.0"
command = f"python {file} --ckpt {sdxl_ckpt}"
run_command(command.split())
command += " --run_compile"
run_command(command.split())
elif file in ["benchmark_controlnet.py", "benchmark_t2i_adapter.py"]:
sdxl_ckpt = (
"diffusers/controlnet-canny-sdxl-1.0"
if "controlnet" in file
else "TencentARC/t2i-adapter-canny-sdxl-1.0"
)
command = f"python {file} --ckpt {sdxl_ckpt}"
run_command(command.split())
command += " --run_compile"
run_command(command.split())
if __name__ == "__main__":
main()
+98
View File
@@ -0,0 +1,98 @@
import argparse
import csv
import gc
import os
from dataclasses import dataclass
from typing import Dict, List, Union
import torch
import torch.utils.benchmark as benchmark
GITHUB_SHA = os.getenv("GITHUB_SHA", None)
BENCHMARK_FIELDS = [
"pipeline_cls",
"ckpt_id",
"batch_size",
"num_inference_steps",
"model_cpu_offload",
"run_compile",
"time (secs)",
"memory (gbs)",
"actual_gpu_memory (gbs)",
"github_sha",
]
PROMPT = "ghibli style, a fantasy landscape with castles"
BASE_PATH = os.getenv("BASE_PATH", ".")
TOTAL_GPU_MEMORY = float(os.getenv("TOTAL_GPU_MEMORY", torch.cuda.get_device_properties(0).total_memory / (1024**3)))
REPO_ID = "diffusers/benchmarks"
FINAL_CSV_FILE = "collated_results.csv"
@dataclass
class BenchmarkInfo:
time: float
memory: float
def flush():
"""Wipes off memory."""
gc.collect()
torch.cuda.empty_cache()
torch.cuda.reset_max_memory_allocated()
torch.cuda.reset_peak_memory_stats()
def bytes_to_giga_bytes(bytes):
return f"{(bytes / 1024 / 1024 / 1024):.3f}"
def benchmark_fn(f, *args, **kwargs):
t0 = benchmark.Timer(
stmt="f(*args, **kwargs)",
globals={"args": args, "kwargs": kwargs, "f": f},
num_threads=torch.get_num_threads(),
)
return f"{(t0.blocked_autorange().mean):.3f}"
def generate_csv_dict(
pipeline_cls: str, ckpt: str, args: argparse.Namespace, benchmark_info: BenchmarkInfo
) -> Dict[str, Union[str, bool, float]]:
"""Packs benchmarking data into a dictionary for latter serialization."""
data_dict = {
"pipeline_cls": pipeline_cls,
"ckpt_id": ckpt,
"batch_size": args.batch_size,
"num_inference_steps": args.num_inference_steps,
"model_cpu_offload": args.model_cpu_offload,
"run_compile": args.run_compile,
"time (secs)": benchmark_info.time,
"memory (gbs)": benchmark_info.memory,
"actual_gpu_memory (gbs)": f"{(TOTAL_GPU_MEMORY):.3f}",
"github_sha": GITHUB_SHA,
}
return data_dict
def write_to_csv(file_name: str, data_dict: Dict[str, Union[str, bool, float]]):
"""Serializes a dictionary into a CSV file."""
with open(file_name, mode="w", newline="") as csvfile:
writer = csv.DictWriter(csvfile, fieldnames=BENCHMARK_FIELDS)
writer.writeheader()
writer.writerow(data_dict)
def collate_csv(input_files: List[str], output_file: str):
"""Collates multiple identically structured CSVs into a single CSV file."""
with open(output_file, mode="w", newline="") as outfile:
writer = csv.DictWriter(outfile, fieldnames=BENCHMARK_FIELDS)
writer.writeheader()
for file in input_files:
with open(file, mode="r") as infile:
reader = csv.DictReader(infile)
for row in reader:
writer.writerow(row)
+15 -1
View File
@@ -72,6 +72,8 @@
title: Overview
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/controlnet
@@ -94,6 +96,8 @@
title: Latent Consistency Model-LoRA
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/svd
title: Stable Video Diffusion
title: Specific pipeline examples
- sections:
- local: training/overview
@@ -129,6 +133,8 @@
title: LoRA
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/lcm_distill
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Methods
@@ -258,6 +264,10 @@
title: ControlNet
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/cycle_diffusion
title: Cycle Diffusion
- local: api/pipelines/dance_diffusion
@@ -278,6 +288,8 @@
title: Kandinsky 2.1
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/kandinsky3
title: Kandinsky 3
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
@@ -327,12 +339,14 @@
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth)
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/adapter
title: Stable Diffusion T2I-Adapter
- local: api/pipelines/stable_diffusion/gligen
+3
View File
@@ -20,6 +20,9 @@ An attention processor is a class for applying different types of attention mech
## AttnProcessor2_0
[[autodoc]] models.attention_processor.AttnProcessor2_0
## FusedAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedAttnProcessor2_0
## LoRAAttnProcessor
[[autodoc]] models.attention_processor.LoRAAttnProcessor
@@ -0,0 +1,39 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
ControlNet-XS generates images with comparable quality to a regular ControlNet, but it is 20-25% faster ([see benchmark](https://github.com/UmerHA/controlnet-xs-benchmark/blob/main/Speed%20Benchmark.ipynb) with StableDiffusion-XL) and uses ~45% less memory.
Here's the overview from the [project page](https://vislearn.github.io/ControlNet-XS/):
*With increasing computing capabilities, current model architectures appear to follow the trend of simply upscaling all components without validating the necessity for doing so. In this project we investigate the size and architectural design of ControlNet [Zhang et al., 2023] for controlling the image generation process with stable diffusion-based models. We show that a new architecture with as little as 1% of the parameters of the base model achieves state-of-the art results, considerably better than ControlNet in terms of FID score. Hence we call it ControlNet-XS. We provide the code for controlling StableDiffusion-XL [Podell et al., 2023] (Model B, 48M Parameters) and StableDiffusion 2.1 [Rombach et al. 2022] (Model B, 14M Parameters), all under openrail license.*
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## StableDiffusionControlNetXSPipeline
[[autodoc]] StableDiffusionControlNetXSPipeline
- all
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
@@ -0,0 +1,45 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS with Stable Diffusion XL
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
ControlNet-XS generates images with comparable quality to a regular ControlNet, but it is 20-25% faster ([see benchmark](https://github.com/UmerHA/controlnet-xs-benchmark/blob/main/Speed%20Benchmark.ipynb)) and uses ~45% less memory.
Here's the overview from the [project page](https://vislearn.github.io/ControlNet-XS/):
*With increasing computing capabilities, current model architectures appear to follow the trend of simply upscaling all components without validating the necessity for doing so. In this project we investigate the size and architectural design of ControlNet [Zhang et al., 2023] for controlling the image generation process with stable diffusion-based models. We show that a new architecture with as little as 1% of the parameters of the base model achieves state-of-the art results, considerably better than ControlNet in terms of FID score. Hence we call it ControlNet-XS. We provide the code for controlling StableDiffusion-XL [Podell et al., 2023] (Model B, 48M Parameters) and StableDiffusion 2.1 [Rombach et al. 2022] (Model B, 14M Parameters), all under openrail license.*
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
<Tip warning={true}>
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) and leave us feedback on how we can improve!
</Tip>
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## StableDiffusionXLControlNetXSPipeline
[[autodoc]] StableDiffusionXLControlNetXSPipeline
- all
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
@@ -0,0 +1,49 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Kandinsky 3
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
The description from it's Github page:
*Kandinsky 3.0 is an open-source text-to-image diffusion model built upon the Kandinsky2-x model family. In comparison to its predecessors, enhancements have been made to the text understanding and visual quality of the model, achieved by increasing the size of the text encoder and Diffusion U-Net models, respectively.*
Its architecture includes 3 main components:
1. [FLAN-UL2](https://huggingface.co/google/flan-ul2), which is an encoder decoder model based on the T5 architecture.
2. New U-Net architecture featuring BigGAN-deep blocks doubles depth while maintaining the same number of parameters.
3. Sber-MoVQGAN is a decoder proven to have superior results in image restoration.
The original codebase can be found at [ai-forever/Kandinsky-3](https://github.com/ai-forever/Kandinsky-3).
<Tip>
Check out the [Kandinsky Community](https://huggingface.co/kandinsky-community) organization on the Hub for the official model checkpoints for tasks like text-to-image, image-to-image, and inpainting.
</Tip>
<Tip>
Make sure to check out the schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Kandinsky3Pipeline
[[autodoc]] Kandinsky3Pipeline
- all
- __call__
## Kandinsky3Img2ImgPipeline
[[autodoc]] Kandinsky3Img2ImgPipeline
- all
- __call__
+5 -1
View File
@@ -40,6 +40,8 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Consistency Models](consistency_models) | unconditional image generation |
| [ControlNet](controlnet) | text2image, image2image, inpainting |
| [ControlNet with Stable Diffusion XL](controlnet_sdxl) | text2image |
| [ControlNet-XS](controlnetxs) | text2image |
| [ControlNet-XS with Stable Diffusion XL](controlnetxs_sdxl) | text2image |
| [Cycle Diffusion](cycle_diffusion) | image2image |
| [Dance Diffusion](dance_diffusion) | unconditional audio generation |
| [DDIM](ddim) | unconditional image generation |
@@ -51,9 +53,10 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
| [Kandinsky 3](kandinsky3) | text2image, image2image |
| [Latent Consistency Models](latent_consistency_models) | text2image |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D, text-to-pano, upscaling |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [Paint by Example](paint_by_example) | inpainting |
@@ -70,6 +73,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Stable Diffusion](stable_diffusion/overview) | text2image, image2image, depth2image, inpainting, image variation, latent upscaler, super-resolution |
| [Stable Diffusion Model Editing](model_editing) | model editing |
| [Stable Diffusion XL](stable_diffusion/stable_diffusion_xl) | text2image, image2image, inpainting |
| [Stable Diffusion XL Turbo](stable_diffusion/sdxl_turbo) | text2image, image2image, inpainting |
| [Stable unCLIP](stable_unclip) | text2image, image variation |
| [Stochastic Karras VE](stochastic_karras_ve) | unconditional image generation |
| [T2I-Adapter](stable_diffusion/adapter) | text2image |
+106
View File
@@ -35,6 +35,112 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
</Tip>
## Inference with under 8GB GPU VRAM
Run the [`PixArtAlphaPipeline`] with under 8GB GPU VRAM by loading the text encoder in 8-bit precision. Let's walk through a full-fledged example.
First, install the [bitsandbytes](https://github.com/TimDettmers/bitsandbytes) library:
```bash
pip install -U bitsandbytes
```
Then load the text encoder in 8-bit:
```python
from transformers import T5EncoderModel
from diffusers import PixArtAlphaPipeline
import torch
text_encoder = T5EncoderModel.from_pretrained(
"PixArt-alpha/PixArt-XL-2-1024-MS",
subfolder="text_encoder",
load_in_8bit=True,
device_map="auto",
)
pipe = PixArtAlphaPipeline.from_pretrained(
"PixArt-alpha/PixArt-XL-2-1024-MS",
text_encoder=text_encoder,
transformer=None,
device_map="auto"
)
```
Now, use the `pipe` to encode a prompt:
```python
with torch.no_grad():
prompt = "cute cat"
prompt_embeds, prompt_attention_mask, negative_embeds, negative_prompt_attention_mask = pipe.encode_prompt(prompt)
```
Since text embeddings have been computed, remove the `text_encoder` and `pipe` from the memory, and free up som GPU VRAM:
```python
import gc
def flush():
gc.collect()
torch.cuda.empty_cache()
del text_encoder
del pipe
flush()
```
Then compute the latents with the prompt embeddings as inputs:
```python
pipe = PixArtAlphaPipeline.from_pretrained(
"PixArt-alpha/PixArt-XL-2-1024-MS",
text_encoder=None,
torch_dtype=torch.float16,
).to("cuda")
latents = pipe(
negative_prompt=None,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_embeds,
prompt_attention_mask=prompt_attention_mask,
negative_prompt_attention_mask=negative_prompt_attention_mask,
num_images_per_prompt=1,
output_type="latent",
).images
del pipe.transformer
flush()
```
<Tip>
Notice that while initializing `pipe`, you're setting `text_encoder` to `None` so that it's not loaded.
</Tip>
Once the latents are computed, pass it off to the VAE to decode into a real image:
```python
with torch.no_grad():
image = pipe.vae.decode(latents / pipe.vae.config.scaling_factor, return_dict=False)[0]
image = pipe.image_processor.postprocess(image, output_type="pil")[0]
image.save("cat.png")
```
By deleting components you aren't using and flushing the GPU VRAM, you should be able to run [`PixArtAlphaPipeline`] with under 8GB GPU VRAM.
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/pixart/8bits_cat.png)
If you want a report of your memory-usage, run this [script](https://gist.github.com/sayakpaul/3ae0f847001d342af27018a96f467e4e).
<Tip warning={true}>
Text embeddings computed in 8-bit can impact the quality of the generated images because of the information loss in the representation space caused by the reduced precision. It's recommended to compare the outputs with and without 8-bit.
</Tip>
While loading the `text_encoder`, you set `load_in_8bit` to `True`. You could also specify `load_in_4bit` to bring your memory requirements down even further to under 7GB.
## PixArtAlphaPipeline
[[autodoc]] PixArtAlphaPipeline
@@ -14,6 +14,11 @@ specific language governing permissions and limitations under the License.
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
Two checkpoints are available for use:
- [ldm3d-original](https://huggingface.co/Intel/ldm3d). The original checkpoint used in the [paper](https://arxiv.org/pdf/2305.10853.pdf)
- [ldm3d-4c](https://huggingface.co/Intel/ldm3d-4c). The new version of LDM3D using 4 channels inputs instead of 6-channels inputs and finetuned on higher resolution images.
The abstract from the paper is:
*This research paper proposes a Latent Diffusion Model for 3D (LDM3D) that generates both image and depth map data from a given text prompt, allowing users to generate RGBD images from text prompts. The LDM3D model is fine-tuned on a dataset of tuples containing an RGB image, depth map and caption, and validated through extensive experiments. We also develop an application called DepthFusion, which uses the generated RGB images and depth maps to create immersive and interactive 360-degree-view experiences using TouchDesigner. This technology has the potential to transform a wide range of industries, from entertainment and gaming to architecture and design. Overall, this paper presents a significant contribution to the field of generative AI and computer vision, and showcases the potential of LDM3D and DepthFusion to revolutionize content creation and digital experiences. A short video summarizing the approach can be found at [this url](https://t.ly/tdi2).*
@@ -26,12 +31,25 @@ Make sure to check out the Stable Diffusion [Tips](overview#tips) section to lea
## StableDiffusionLDM3DPipeline
[[autodoc]] StableDiffusionLDM3DPipeline
[[autodoc]] pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d.StableDiffusionLDM3DPipeline
- all
- __call__
## LDM3DPipelineOutput
[[autodoc]] pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d.LDM3DPipelineOutput
- all
- __call__
# Upscaler
[LDM3D-VR](https://arxiv.org/pdf/2311.03226.pdf) is an extended version of LDM3D.
The abstract from the paper is:
*Latent diffusion models have proven to be state-of-the-art in the creation and manipulation of visual outputs. However, as far as we know, the generation of depth maps jointly with RGB is still limited. We introduce LDM3D-VR, a suite of diffusion models targeting virtual reality development that includes LDM3D-pano and LDM3D-SR. These models enable the generation of panoramic RGBD based on textual prompts and the upscaling of low-resolution inputs to high-resolution RGBD, respectively. Our models are fine-tuned from existing pretrained models on datasets containing panoramic/high-resolution RGB images, depth maps and captions. Both models are evaluated in comparison to existing related methods*
Two checkpoints are available for use:
- [ldm3d-pano](https://huggingface.co/Intel/ldm3d-pano). This checkpoint enables the generation of panoramic images and requires the StableDiffusionLDM3DPipeline pipeline to be used.
- [ldm3d-sr](https://huggingface.co/Intel/ldm3d-sr). This checkpoint enables the upscaling of RGB and depth images. Can be used in cascade after the original LDM3D pipeline using the StableDiffusionUpscaleLDM3DPipeline from communauty pipeline.
@@ -121,10 +121,16 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2 text-gray-700">
<a href="./ldm3d_diffusion">StableDiffusionLDM3D</a>
</td>
<td class="px-4 py-2 text-gray-700">text-to-rgb, text-to-depth</td>
<td class="px-4 py-2 text-gray-700">text-to-rgb, text-to-depth, text-to-pano</td>
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/r23/ldm3d-space"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./ldm3d_diffusion">StableDiffusionUpscaleLDM3D</a>
</td>
<td class="px-4 py-2 text-gray-700">ldm3d super-resolution</td>
</tr>
</tbody>
</table>
</div>
@@ -0,0 +1,35 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# SDXL Turbo
Stable Diffusion XL (SDXL) Turbo was proposed in [Adversarial Diffusion Distillation](https://stability.ai/research/adversarial-diffusion-distillation) by Axel Sauer, Dominik Lorenz, Andreas Blattmann, and Robin Rombach.
The abstract from the paper is:
*We introduce Adversarial Diffusion Distillation (ADD), a novel training approach that efficiently samples large-scale foundational image diffusion models in just 14 steps while maintaining high image quality. We use score distillation to leverage large-scale off-the-shelf image diffusion models as a teacher signal in combination with an adversarial loss to ensure high image fidelity even in the low-step regime of one or two sampling steps. Our analyses show that our model clearly outperforms existing few-step methods (GANs,Latent Consistency Models) in a single step and reaches the performance of state-of-the-art diffusion models (SDXL) in only four steps. ADD is the first method to unlock single-step, real-time image synthesis with foundation models.*
## Tips
- SDXL Turbo uses the exact same architecture as [SDXL](./stable_diffusion_xl), which means it also has the same API. Please refer to the [SDXL](./stable_diffusion_xl) API reference for more details.
- SDXL Turbo should disable guidance scale by setting `guidance_scale=0.0`
- SDXL Turbo should use `timestep_spacing='trailing'` for the scheduler and use between 1 and 4 steps.
- SDXL Turbo has been trained to generate images of size 512x512.
- SDXL Turbo is open-access, but not open-source meaning that one might have to buy a model license in order to use it for commercial applications. Make sure to read the [official model card](https://huggingface.co/stabilityai/sdxl-turbo) to learn more.
<Tip>
To learn how to use SDXL Turbo for various tasks, how to optimize performance, and other usage examples, take a look at the [SDXL Turbo](../../../using-diffusers/sdxl_turbo) guide.
Check out the [Stability AI](https://huggingface.co/stabilityai) Hub organization for the official base and refiner model checkpoints!
</Tip>
@@ -92,6 +92,19 @@ imageio.mimsave("video.mp4", result, fps=4)
```
- #### SDXL Support
In order to use the SDXL model when generating a video from prompt, use the `TextToVideoZeroSDXLPipeline` pipeline:
```python
import torch
from diffusers import TextToVideoZeroSDXLPipeline
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
pipe = TextToVideoZeroSDXLPipeline.from_pretrained(
model_id, torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
```
### Text-To-Video with Pose Control
To generate a video from prompt with additional pose control
@@ -141,7 +154,33 @@ To generate a video from prompt with additional pose control
result = pipe(prompt=[prompt] * len(pose_images), image=pose_images, latents=latents).images
imageio.mimsave("video.mp4", result, fps=4)
```
- #### SDXL Support
Since our attention processor also works with SDXL, it can be utilized to generate a video from prompt using ControlNet models powered by SDXL:
```python
import torch
from diffusers import StableDiffusionXLControlNetPipeline, ControlNetModel
from diffusers.pipelines.text_to_video_synthesis.pipeline_text_to_video_zero import CrossFrameAttnProcessor
controlnet_model_id = 'thibaud/controlnet-openpose-sdxl-1.0'
model_id = 'stabilityai/stable-diffusion-xl-base-1.0'
controlnet = ControlNetModel.from_pretrained(controlnet_model_id, torch_dtype=torch.float16)
pipe = StableDiffusionControlNetPipeline.from_pretrained(
model_id, controlnet=controlnet, torch_dtype=torch.float16
).to('cuda')
# Set the attention processor
pipe.unet.set_attn_processor(CrossFrameAttnProcessor(batch_size=2))
pipe.controlnet.set_attn_processor(CrossFrameAttnProcessor(batch_size=2))
# fix latents for all frames
latents = torch.randn((1, 4, 128, 128), device="cuda", dtype=torch.float16).repeat(len(pose_images), 1, 1, 1)
prompt = "Darth Vader dancing in a desert"
result = pipe(prompt=[prompt] * len(pose_images), image=pose_images, latents=latents).images
imageio.mimsave("video.mp4", result, fps=4)
```
### Text-To-Video with Edge Control
@@ -253,5 +292,10 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## TextToVideoZeroSDXLPipeline
[[autodoc]] TextToVideoZeroSDXLPipeline
- all
- __call__
## TextToVideoPipelineOutput
[[autodoc]] pipelines.text_to_video_synthesis.pipeline_text_to_video_zero.TextToVideoPipelineOutput
@@ -24,7 +24,7 @@ The abstract from the paper is:
*Model-based reinforcement learning methods often use learning only for the purpose of estimating an approximate dynamics model, offloading the rest of the decision-making work to classical trajectory optimizers. While conceptually simple, this combination has a number of empirical shortcomings, suggesting that learned models may not be well-suited to standard trajectory optimization. In this paper, we consider what it would look like to fold as much of the trajectory optimization pipeline as possible into the modeling problem, such that sampling from the model and planning with it become nearly identical. The core of our technical approach lies in a diffusion probabilistic model that plans by iteratively denoising trajectories. We show how classifier-guided sampling and image inpainting can be reinterpreted as coherent planning strategies, explore the unusual and useful properties of diffusion-based planning methods, and demonstrate the effectiveness of our framework in control settings that emphasize long-horizon decision-making and test-time flexibility.*
You can find additional information about the model on the [project page](https://diffusion-planning.github.io/), the [original codebase](https://github.com/jannerm/diffuser), or try it out in a demo [notebook](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb).
You can find additional information about the model on the [project page](https://diffusion-planning.github.io/), the [original codebase](https://github.com/jannerm/diffuser), or try it out in a demo [notebook](https://colab.research.google.com/drive/1rXm8CX4ZdN5qivjJ2lhwhkOmt_m0CvU0#scrollTo=6HXJvhyqcITc&uniqifier=1).
The script to run the model is available [here](https://github.com/huggingface/diffusers/tree/main/examples/reinforcement_learning).
@@ -25,4 +25,4 @@ The abstract from the paper is:
</Tip>
## ScoreSdeVpScheduler
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler
[[autodoc]] schedulers.deprecated.scheduling_sde_vp.ScoreSdeVpScheduler
@@ -18,4 +18,4 @@ specific language governing permissions and limitations under the License.
[[autodoc]] KarrasVeScheduler
## KarrasVeOutput
[[autodoc]] schedulers.scheduling_karras_ve.KarrasVeOutput
[[autodoc]] schedulers.deprecated.scheduling_karras_ve.KarrasVeOutput
+27 -7
View File
@@ -297,17 +297,37 @@ if you don't know yet what specific component you would like to add:
- [Model or pipeline](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+pipeline%2Fmodel%22)
- [Scheduler](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+scheduler%22)
Before adding any of the three components, it is strongly recommended that you give the [Philosophy guide](philosophy) a read to better understand the design of any of the three components. Please be aware that
we cannot merge model, scheduler, or pipeline additions that strongly diverge from our design philosophy
as it will lead to API inconsistencies. If you fundamentally disagree with a design choice, please
open a [Feedback issue](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=) instead so that it can be discussed whether a certain design
pattern/design choice shall be changed everywhere in the library and whether we shall update our design philosophy. Consistency across the library is very important for us.
Before adding any of the three components, it is strongly recommended that you give the [Philosophy guide](philosophy) a read to better understand the design of any of the three components. Please be aware that we cannot merge model, scheduler, or pipeline additions that strongly diverge from our design philosophy
as it will lead to API inconsistencies. If you fundamentally disagree with a design choice, please open a [Feedback issue](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=) instead so that it can be discussed whether a certain design pattern/design choice shall be changed everywhere in the library and whether we shall update our design philosophy. Consistency across the library is very important for us.
Please make sure to add links to the original codebase/paper to the PR and ideally also ping the
original author directly on the PR so that they can follow the progress and potentially help with questions.
Please make sure to add links to the original codebase/paper to the PR and ideally also ping the original author directly on the PR so that they can follow the progress and potentially help with questions.
If you are unsure or stuck in the PR, don't hesitate to leave a message to ask for a first review or help.
#### Copied from mechanism
A unique and important feature to understand when adding any pipeline, model or scheduler code is the `# Copied from` mechanism. You'll see this all over the Diffusers codebase, and the reason we use it is to keep the codebase easy to understand and maintain. Marking code with the `# Copied from` mechanism forces the marked code to be identical to the code it was copied from. This makes it easy to update and propagate changes across many files whenever you run `make fix-copies`.
For example, in the code example below, [`~diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput`] is the original code and `AltDiffusionPipelineOutput` uses the `# Copied from` mechanism to copy it. The only difference is changing the class prefix from `Stable` to `Alt`.
```py
# Copied from diffusers.pipelines.stable_diffusion.pipeline_output.StableDiffusionPipelineOutput with Stable->Alt
class AltDiffusionPipelineOutput(BaseOutput):
"""
Output class for Alt Diffusion pipelines.
Args:
images (`List[PIL.Image.Image]` or `np.ndarray`)
List of denoised PIL images of length `batch_size` or NumPy array of shape `(batch_size, height, width,
num_channels)`.
nsfw_content_detected (`List[bool]`)
List indicating whether the corresponding generated image contains "not-safe-for-work" (nsfw) content or
`None` if safety checking could not be performed.
"""
```
To learn more, read this section of the [~Don't~ Repeat Yourself*](https://huggingface.co/blog/transformers-design-philosophy#4-machine-learning-models-are-static) blog post.
## How to write a good issue
**The better your issue is written, the higher the chances that it will be quickly resolved.**
+255
View File
@@ -0,0 +1,255 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Latent Consistency Distillation
[Latent Consistency Models (LCMs)](https://hf.co/papers/2310.04378) are able to generate high-quality images in just a few steps, representing a big leap forward because many pipelines require at least 25+ steps. LCMs are produced by applying the latent consistency distillation method to any Stable Diffusion model. This method works by applying *one-stage guided distillation* to the latent space, and incorporating a *skipping-step* method to consistently skip timesteps to accelerate the distillation process (refer to section 4.1, 4.2, and 4.3 of the paper for more details).
If you're training on a GPU with limited vRAM, try enabling `gradient_checkpointing`, `gradient_accumulation_steps`, and `mixed_precision` to reduce memory-usage and speedup training. You can reduce your memory-usage even more by enabling memory-efficient attention with [xFormers](../optimization/xformers) and [bitsandbytes'](https://github.com/TimDettmers/bitsandbytes) 8-bit optimizer.
This guide will explore the [train_lcm_distill_sd_wds.py](https://github.com/huggingface/diffusers/blob/main/examples/consistency_distillation/train_lcm_distill_sd_wds.py) script to help you become more familiar with it, and how you can adapt it for your own use-case.
Before running the script, make sure you install the library from source:
```bash
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install .
```
Then navigate to the example folder containing the training script and install the required dependencies for the script you're using:
```bash
cd examples/consistency_distillation
pip install -r requirements.txt
```
<Tip>
🤗 Accelerate is a library for helping you train on multiple GPUs/TPUs or with mixed-precision. It'll automatically configure your training setup based on your hardware and environment. Take a look at the 🤗 Accelerate [Quick tour](https://huggingface.co/docs/accelerate/quicktour) to learn more.
</Tip>
Initialize an 🤗 Accelerate environment (try enabling `torch.compile` to significantly speedup training):
```bash
accelerate config
```
To setup a default 🤗 Accelerate environment without choosing any configurations:
```bash
accelerate config default
```
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```bash
from accelerate.utils import write_basic_config
write_basic_config()
```
Lastly, if you want to train a model on your own dataset, take a look at the [Create a dataset for training](create_dataset) guide to learn how to create a dataset that works with the training script.
## Script parameters
<Tip>
The following sections highlight parts of the training script that are important for understanding how to modify it, but it doesn't cover every aspect of the script in detail. If you're interested in learning more, feel free to read through the [script](https://github.com/huggingface/diffusers/blob/main/examples/consistency_distillation/train_lcm_distill_sd_wds.py) and let us know if you have any questions or concerns.
</Tip>
The training script provides many parameters to help you customize your training run. All of the parameters and their descriptions are found in the [`parse_args()`](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L419) function. This function provides default values for each parameter, such as the training batch size and learning rate, but you can also set your own values in the training command if you'd like.
For example, to speedup training with mixed precision using the fp16 format, add the `--mixed_precision` parameter to the training command:
```bash
accelerate launch train_lcm_distill_sd_wds.py \
--mixed_precision="fp16"
```
Most of the parameters are identical to the parameters in the [Text-to-image](text2image#script-parameters) training guide, so you'll focus on the parameters that are relevant to latent consistency distillation in this guide.
- `--pretrained_teacher_model`: the path to a pretrained latent diffusion model to use as the teacher model
- `--pretrained_vae_model_name_or_path`: path to a pretrained VAE; the SDXL VAE is known to suffer from numerical instability, so this parameter allows you to specify an alternative VAE (like this [VAE]((https://huggingface.co/madebyollin/sdxl-vae-fp16-fix)) by madebyollin which works in fp16)
- `--w_min` and `--w_max`: the minimum and maximum guidance scale values for guidance scale sampling
- `--num_ddim_timesteps`: the number of timesteps for DDIM sampling
- `--loss_type`: the type of loss (L2 or Huber) to calculate for latent consistency distillation; Huber loss is generally preferred because it's more robust to outliers
- `--huber_c`: the Huber loss parameter
## Training script
The training script starts by creating a dataset class - [`Text2ImageDataset`](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L141) - for preprocessing the images and creating a training dataset.
```py
def transform(example):
image = example["image"]
image = TF.resize(image, resolution, interpolation=transforms.InterpolationMode.BILINEAR)
c_top, c_left, _, _ = transforms.RandomCrop.get_params(image, output_size=(resolution, resolution))
image = TF.crop(image, c_top, c_left, resolution, resolution)
image = TF.to_tensor(image)
image = TF.normalize(image, [0.5], [0.5])
example["image"] = image
return example
```
For improved performance on reading and writing large datasets stored in the cloud, this script uses the [WebDataset](https://github.com/webdataset/webdataset) format to create a preprocessing pipeline to apply transforms and create a dataset and dataloader for training. Images are processed and fed to the training loop without having to download the full dataset first.
```py
processing_pipeline = [
wds.decode("pil", handler=wds.ignore_and_continue),
wds.rename(image="jpg;png;jpeg;webp", text="text;txt;caption", handler=wds.warn_and_continue),
wds.map(filter_keys({"image", "text"})),
wds.map(transform),
wds.to_tuple("image", "text"),
]
```
In the [`main()`](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L768) function, all the necessary components like the noise scheduler, tokenizers, text encoders, and VAE are loaded. The teacher UNet is also loaded here and then you can create a student UNet from the teacher UNet. The student UNet is updated by the optimizer during training.
```py
teacher_unet = UNet2DConditionModel.from_pretrained(
args.pretrained_teacher_model, subfolder="unet", revision=args.teacher_revision
)
unet = UNet2DConditionModel(**teacher_unet.config)
unet.load_state_dict(teacher_unet.state_dict(), strict=False)
unet.train()
```
Now you can create the [optimizer](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L979) to update the UNet parameters:
```py
optimizer = optimizer_class(
unet.parameters(),
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
)
```
Create the [dataset](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L994):
```py
dataset = Text2ImageDataset(
train_shards_path_or_url=args.train_shards_path_or_url,
num_train_examples=args.max_train_samples,
per_gpu_batch_size=args.train_batch_size,
global_batch_size=args.train_batch_size * accelerator.num_processes,
num_workers=args.dataloader_num_workers,
resolution=args.resolution,
shuffle_buffer_size=1000,
pin_memory=True,
persistent_workers=True,
)
train_dataloader = dataset.train_dataloader
```
Next, you're ready to setup the [training loop](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L1049) and implement the latent consistency distillation method (see Algorithm 1 in the paper for more details). This section of the script takes care of adding noise to the latents, sampling and creating a guidance scale embedding, and predicting the original image from the noise.
```py
pred_x_0 = predicted_origin(
noise_pred,
start_timesteps,
noisy_model_input,
noise_scheduler.config.prediction_type,
alpha_schedule,
sigma_schedule,
)
model_pred = c_skip_start * noisy_model_input + c_out_start * pred_x_0
```
It gets the [teacher model predictions](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L1172) and the [LCM predictions](https://github.com/huggingface/diffusers/blob/3b37488fa3280aed6a95de044d7a42ffdcb565ef/examples/consistency_distillation/train_lcm_distill_sd_wds.py#L1209) next, calculates the loss, and then backpropagates it to the LCM.
```py
if args.loss_type == "l2":
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
elif args.loss_type == "huber":
loss = torch.mean(
torch.sqrt((model_pred.float() - target.float()) ** 2 + args.huber_c**2) - args.huber_c
)
```
If you want to learn more about how the training loop works, check out the [Understanding pipelines, models and schedulers tutorial](../using-diffusers/write_own_pipeline) which breaks down the basic pattern of the denoising process.
## Launch the script
Now you're ready to launch the training script and start distilling!
For this guide, you'll use the `--train_shards_path_or_url` to specify the path to the [Conceptual Captions 12M](https://github.com/google-research-datasets/conceptual-12m) dataset stored on the Hub [here](https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset). Set the `MODEL_DIR` environment variable to the name of the teacher model and `OUTPUT_DIR` to where you want to save the model.
```bash
export MODEL_DIR="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_sd_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
--resolution=512 \
--learning_rate=1e-6 --loss_type="huber" --ema_decay=0.95 --adam_weight_decay=0.0 \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
--gradient_checkpointing --enable_xformers_memory_efficient_attention \
--gradient_accumulation_steps=1 \
--use_8bit_adam \
--resume_from_checkpoint=latest \
--report_to=wandb \
--seed=453645634 \
--push_to_hub
```
Once training is complete, you can use your new LCM for inference.
```py
from diffusers import UNet2DConditionModel, DiffusionPipeline, LCMScheduler
import torch
unet = UNet2DConditionModel.from_pretrained("your-username/your-model", torch_dtype=torch.float16, variant="fp16")
pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", unet=unet, torch_dtype=torch.float16, variant="fp16")
pipeline.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
pipeline.to("cuda")
prompt = "sushi rolls in the form of panda heads, sushi platter"
image = pipeline(prompt, num_inference_steps=4, guidance_scale=1.0).images[0]
```
## LoRA
LoRA is a training technique for significantly reducing the number of trainable parameters. As a result, training is faster and it is easier to store the resulting weights because they are a lot smaller (~100MBs). Use the [train_lcm_distill_lora_sd_wds.py](https://github.com/huggingface/diffusers/blob/main/examples/consistency_distillation/train_lcm_distill_lora_sd_wds.py) or [train_lcm_distill_lora_sdxl.wds.py](https://github.com/huggingface/diffusers/blob/main/examples/consistency_distillation/train_lcm_distill_lora_sdxl_wds.py) script to train with LoRA.
The LoRA training script is discussed in more detail in the [LoRA training](lora) guide.
## Stable Diffusion XL
Stable Diffusion XL (SDXL) is a powerful text-to-image model that generates high-resolution images, and it adds a second text-encoder to its architecture. Use the [train_lcm_distill_sdxl_wds.py](https://github.com/huggingface/diffusers/blob/main/examples/consistency_distillation/train_lcm_distill_sdxl_wds.py) script to train a SDXL model with LoRA.
The SDXL training script is discussed in more detail in the [SDXL training](sdxl) guide.
## Next steps
Congratulations on distilling a LCM model! To learn more about LCM, the following may be helpful:
- Learn how to use [LCMs for inference](../using-diffusers/lcm) for text-to-image, image-to-image, and with LoRA checkpoints.
- Read the [SDXL in 4 steps with Latent Consistency LoRAs](https://huggingface.co/blog/lcm_lora) blog post to learn more about SDXL LCM-LoRA's for super fast inference, quality comparisons, benchmarks, and more.
@@ -20,6 +20,8 @@ The Kandinsky models are a series of multilingual text-to-image generation model
[Kandinsky 2.2](../api/pipelines/kandinsky_v22) improves on the previous model by replacing the image encoder of the image prior model with a larger CLIP-ViT-G model to improve quality. The image prior model was also retrained on images with different resolutions and aspect ratios to generate higher-resolution images and different image sizes.
[Kandinsky 3](../api/pipelines/kandinsky3) simplifies the architecture and shifts away from the two-stage generation process involving the prior model and diffusion model. Instead, Kandinsky 3 uses [Flan-UL2](https://huggingface.co/google/flan-ul2) to encode text, a UNet with [BigGan-deep](https://hf.co/papers/1809.11096) blocks, and [Sber-MoVQGAN](https://github.com/ai-forever/MoVQGAN) to decode the latents into images. Text understanding and generated image quality are primarily achieved by using a larger text encoder and UNet.
This guide will show you how to use the Kandinsky models for text-to-image, image-to-image, inpainting, interpolation, and more.
Before you begin, make sure you have the following libraries installed:
@@ -33,6 +35,10 @@ Before you begin, make sure you have the following libraries installed:
Kandinsky 2.1 and 2.2 usage is very similar! The only difference is Kandinsky 2.2 doesn't accept `prompt` as an input when decoding the latents. Instead, Kandinsky 2.2 only accepts `image_embeds` during decoding.
<br>
Kandinsky 3 has a more concise architecture and it doesn't require a prior model. This means it's usage is identical to other diffusion models like [Stable Diffusion XL](sdxl).
</Tip>
## Text-to-image
@@ -91,6 +97,23 @@ image
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-text-to-image.png"/>
</div>
</hfoption>
<hfoption id="Kandinsky 3">
Kandinsky 3 doesn't require a prior model so you can directly load the [`Kandinsky3Pipeline`] and pass a prompt to generate an image:
```py
from diffusers import Kandinsky3Pipeline
import torch
pipeline = Kandinsky3Pipeline.from_pretrained("kandinsky-community/kandinsky-3", variant="fp16", torch_dtype=torch.float16)
pipeline.enable_model_cpu_offload()
prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
image = pipeline(prompt).images[0]
image
```
</hfoption>
</hfoptions>
@@ -161,6 +184,20 @@ prior_pipeline = KandinskyPriorPipeline.from_pretrained("kandinsky-community/kan
pipeline = KandinskyV22Img2ImgPipeline.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16, use_safetensors=True).to("cuda")
```
</hfoption>
<hfoption id="Kandinsky 3">
Kandinsky 3 doesn't require a prior model so you can directly load the image-to-image pipeline:
```py
from diffusers import Kandinsky3Img2ImgPipeline
from diffusers.utils import load_image
import torch
pipeline = Kandinsky3Img2ImgPipeline.from_pretrained("kandinsky-community/kandinsky-3", variant="fp16", torch_dtype=torch.float16)
pipeline.enable_model_cpu_offload()
```
</hfoption>
</hfoptions>
@@ -218,6 +255,14 @@ make_image_grid([original_image.resize((512, 512)), image.resize((512, 512))], r
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-image-to-image.png"/>
</div>
</hfoption>
<hfoption id="Kandinsky 3">
```py
image = pipeline(prompt, negative_prompt=negative_prompt, image=image, strength=0.75, num_inference_steps=25).images[0]
image
```
</hfoption>
</hfoptions>
@@ -307,3 +307,394 @@ prompt = "a house by william eggleston, sunrays, beautiful, sunlight, sunrays, b
image = pipeline(prompt=prompt).images[0]
image
```
## IP-Adapter
[IP-Adapter](https://ip-adapter.github.io/) is an effective and lightweight adapter that adds image prompting capabilities to a diffusion model. This adapter works by decoupling the cross-attention layers of the image and text features. All the other model components are frozen and only the embedded image features in the UNet are trained. As a result, IP-Adapter files are typically only ~100MBs.
IP-Adapter works with most of our pipelines, including Stable Diffusion, Stable Diffusion XL (SDXL), ControlNet, T2I-Adapter, AnimateDiff. And you can use any custom models finetuned from the same base models. It also works with LCM-Lora out of box.
<Tip>
You can find official IP-Adapter checkpoints in [h94/IP-Adapter](https://huggingface.co/h94/IP-Adapter).
IP-Adapter was contributed by [okotaku](https://github.com/okotaku).
</Tip>
Let's first create a Stable Diffusion Pipeline.
```py
from diffusers import AutoPipelineForText2Image
import torch
from diffusers.utils import load_image
pipeline = AutoPipelineForText2Image.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to("cuda")
```
Now load the [h94/IP-Adapter](https://huggingface.co/h94/IP-Adapter) weights with the [`~loaders.IPAdapterMixin.load_ip_adapter`] method.
```py
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
```
<Tip>
IP-Adapter relies on an image encoder to generate the image features, if your IP-Adapter weights folder contains a "image_encoder" subfolder, the image encoder will be automatically loaded and registered to the pipeline. Otherwise you can so load a [`~transformers.CLIPVisionModelWithProjection`] model and pass it to a Stable Diffusion pipeline when you create it.
```py
from diffusers import AutoPipelineForText2Image, CLIPVisionModelWithProjection
import torch
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"h94/IP-Adapter",
subfolder="models/image_encoder",
torch_dtype=torch.float16,
).to("cuda")
pipeline = AutoPipelineForText2Image.from_pretrained("runwayml/stable-diffusion-v1-5", image_encoder=image_encoder, torch_dtype=torch.float16).to("cuda")
```
</Tip>
IP-Adapter allows you to use both image and text to condition the image generation process. For example, let's use the bear image from the [Textual Inversion](#textual-inversion) section as the image prompt (`ip_adapter_image`) along with a text prompt to add "sunglasses". 😎
```py
pipeline.set_ip_adapter_scale(0.6)
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_neg_embed.png")
generator = torch.Generator(device="cpu").manual_seed(33)
images = pipeline(
    prompt='best quality, high quality, wearing sunglasses',
    ip_adapter_image=image,
    negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
    num_inference_steps=50,
    generator=generator,
).images
images[0]
```
<div class="flex justify-center">
    <img src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip-bear.png" />
</div>
<Tip>
You can use the [`~loaders.IPAdapterMixin.set_ip_adapter_scale`] method to adjust the text prompt and image prompt condition ratio.  If you're only using the image prompt, you should set the scale to `1.0`. You can lower the scale to get more generation diversity, but it'll be less aligned with the prompt.
`scale=0.5` can achieve good results in most cases when you use both text and image prompts.
</Tip>
IP-Adapter also works great with Image-to-Image and Inpainting pipelines. See below examples of how you can use it with Image-to-Image and Inpaint.
<hfoptions id="tasks">
<hfoption id="image-to-image">
```py
from diffusers import AutoPipelineForImage2Image
import torch
from diffusers.utils import load_image
pipeline = AutoPipelineForImage2Image.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to("cuda")
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/vermeer.jpg")
ip_image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/river.png")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
generator = torch.Generator(device="cpu").manual_seed(33)
images = pipeline(
    prompt='best quality, high quality',
    image = image,
    ip_adapter_image=ip_image,
    num_inference_steps=50,
    generator=generator,
    strength=0.6,
).images
images[0]
```
</hfoption>
<hfoption id="inpaint">
```py
from diffusers import AutoPipelineForInpaint
import torch
from diffusers.utils import load_image
pipeline = AutoPipelineForInpaint.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float).to("cuda")
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/inpaint_image.png")
mask = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/mask.png")
ip_image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/girl.png")
image = image.resize((512, 768))
mask = mask.resize((512, 768))
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
generator = torch.Generator(device="cpu").manual_seed(33)
images = pipeline(
prompt='best quality, high quality',
image = image,
mask_image = mask,
ip_adapter_image=ip_image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50,
generator=generator,
strength=0.5,
).images
images[0]
```
</hfoption>
</hfoptions>
IP-Adapters can also be used with [SDXL](../api/pipelines/stable_diffusion/stable_diffusion_xl.md)
```python
from diffusers import AutoPipelineForText2Image
from diffusers.utils import load_image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16
).to("cuda")
image = load_image("https://huggingface.co/datasets/sayakpaul/sample-datasets/resolve/main/watercolor_painting.jpeg")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name="ip-adapter_sdxl.bin")
generator = torch.Generator(device="cpu").manual_seed(33)
image = pipeline(
prompt="best quality, high quality",
ip_adapter_image=image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=25,
generator=generator,
).images[0]
image.save("sdxl_t2i.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/sayakpaul/sample-datasets/resolve/main/watercolor_painting.jpeg"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/sayakpaul/sample-datasets/resolve/main/sdxl_t2i.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">adapted image</figcaption>
</div>
</div>
You can use the IP-Adapter face model to apply specific faces to your images. It is an effective way to maintain consistent characters in your image generations.
Weights are loaded with the same method used for the other IP-Adapters.
```python
# Load ip-adapter-full-face_sd15.bin
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter-full-face_sd15.bin")
```
<Tip>
It is recommended to use `DDIMScheduler` and `EulerDiscreteScheduler` for face model.
</Tip>
```python
import torch
from diffusers import StableDiffusionPipeline, DDIMScheduler
from diffusers.utils import load_image
noise_scheduler = DDIMScheduler(
num_train_timesteps=1000,
beta_start=0.00085,
beta_end=0.012,
beta_schedule="scaled_linear",
clip_sample=False,
set_alpha_to_one=False,
steps_offset=1
)
pipeline = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
scheduler=noise_scheduler,
).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter-full-face_sd15.bin")
pipeline.set_ip_adapter_scale(0.7)
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png")
generator = torch.Generator(device="cpu").manual_seed(33)
image = pipeline(
prompt="A photo of a girl wearing a black dress, holding red roses in hand, upper body, behind is the Eiffel Tower",
ip_adapter_image=image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50, num_images_per_prompt=1, width=512, height=704,
generator=generator,
).images[0]
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ai_face2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ipadapter_full_face_output.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">output image</figcaption>
</div>
</div>
### LCM-Lora
You can use IP-Adapter with LCM-Lora to achieve "instant fine-tune" with custom images. Note that you need to load IP-Adapter weights before loading the LCM-Lora weights.
```py
from diffusers import DiffusionPipeline, LCMScheduler
import torch
from diffusers.utils import load_image
model_id = "sd-dreambooth-library/herge-style"
lcm_lora_id = "latent-consistency/lcm-lora-sdv1-5"
pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipe.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
pipe.load_lora_weights(lcm_lora_id)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
pipe.enable_model_cpu_offload()
prompt = "best quality, high quality"
image = load_image("https://user-images.githubusercontent.com/24734142/266492875-2d50d223-8475-44f0-a7c6-08b51cb53572.png")
images = pipe(
prompt=prompt,
ip_adapter_image=image,
num_inference_steps=4,
guidance_scale=1,
).images[0]
```
### Other pipelines
IP-Adapter is compatible with any pipeline that (1) uses a text prompt and (2) uses Stable Diffusion or Stable Diffusion XL checkpoint. To use IP-Adapter with a different pipeline, all you need to do is to run `load_ip_adapter()` method after you create the pipeline, and then pass your image to the pipeline as `ip_adapter_image`
<Tip>
🤗 Diffusers currently only supports using IP-Adapter with some of the most popular pipelines, feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you have a cool use-case and require integrating IP-adapters with a pipeline that does not support it yet!
</Tip>
You can find below examples on how to use IP-Adapter with ControlNet and AnimateDiff.
<hfoptions id="model">
<hfoption id="ControlNet">
```
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
import torch
from diffusers.utils import load_image
controlnet_model_path = "lllyasviel/control_v11f1p_sd15_depth"
controlnet = ControlNetModel.from_pretrained(controlnet_model_path, torch_dtype=torch.float16)
pipeline = StableDiffusionControlNetPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16)
pipeline.to("cuda")
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/statue.png")
depth_map = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/depth.png")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
generator = torch.Generator(device="cpu").manual_seed(33)
images = pipeline(
prompt='best quality, high quality',
image=depth_map,
ip_adapter_image=image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50,
generator=generator,
).images
images[0]
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/statue.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ipa-controlnet-out.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">adapted image</figcaption>
</div>
</div>
</hfoption>
<hfoption id="AnimateDiff">
```py
# animate diff + ip adapter
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from diffusers.utils import export_to_gif, load_image
# Load the motion adapter
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2", torch_dtype=torch.float16)
# load SD 1.5 based finetuned model
model_id = "Lykon/DreamShaper"
pipe = AnimateDiffPipeline.from_pretrained(model_id, motion_adapter=adapter, torch_dtype=torch.float16)
# scheduler
scheduler = DDIMScheduler(
clip_sample=False,
beta_start=0.00085,
beta_end=0.012,
beta_schedule="linear",
timestep_spacing="trailing",
steps_offset=1
)
pipe.scheduler = scheduler
# enable memory savings
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
# load ip_adapter
pipe.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
# load motion adapters
pipe.load_lora_weights("guoyww/animatediff-motion-lora-zoom-out", adapter_name="zoom-out")
pipe.load_lora_weights("guoyww/animatediff-motion-lora-tilt-up", adapter_name="tilt-up")
pipe.load_lora_weights("guoyww/animatediff-motion-lora-pan-left", adapter_name="pan-left")
seed = 42
image = load_image("https://user-images.githubusercontent.com/24734142/266492875-2d50d223-8475-44f0-a7c6-08b51cb53572.png")
images = [image] * 3
prompts = ["best quality, high quality"] * 3
negative_prompt = "bad quality, worst quality"
adapter_weights = [[0.75, 0.0, 0.0], [0.0, 0.0, 0.75], [0.0, 0.75, 0.75]]
# generate
output_frames = []
for prompt, image, adapter_weight in zip(prompts, images, adapter_weights):
pipe.set_adapters(["zoom-out", "tilt-up", "pan-left"], adapter_weights=adapter_weight)
output = pipe(
prompt= prompt,
num_frames=16,
guidance_scale=7.5,
num_inference_steps=30,
ip_adapter_image = image,
generator=torch.Generator("cpu").manual_seed(seed),
)
frames = output.frames[0]
output_frames.extend(frames)
export_to_gif(output_frames, "test_out_animation.gif")
```
</hfoption>
</hfoptions>
@@ -174,10 +174,4 @@ Set `private=True` in the [`~diffusers.utils.PushToHubMixin.push_to_hub`] functi
controlnet.push_to_hub("my-controlnet-model-private", private=True)
```
Private repositories are only visible to you, and other users won't be able to clone the repository and your repository won't appear in search results. Even if a user has the URL to your private repository, they'll receive a `404 - Sorry, we can't find the page you are looking for.`
To load a model, scheduler, or pipeline from private or gated repositories, set `use_auth_token=True`:
```py
model = ControlNetModel.from_pretrained("your-namespace/my-controlnet-model-private", use_auth_token=True)
```
Private repositories are only visible to you, and other users won't be able to clone the repository and your repository won't appear in search results. Even if a user has the URL to your private repository, they'll receive a `404 - Sorry, we can't find the page you are looking for`. You must be [logged in](https://huggingface.co/docs/huggingface_hub/quick-start#login) to load a model from a private repository.
@@ -0,0 +1,116 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Diffusion XL Turbo
[[open-in-colab]]
SDXL Turbo is an adversarial time-distilled [Stable Diffusion XL](https://huggingface.co/papers/2307.01952) (SDXL) model capable
of running inference in as little as 1 step.
This guide will show you how to use SDXL-Turbo for text-to-image and image-to-image.
Before you begin, make sure you have the following libraries installed:
```py
# uncomment to install the necessary libraries in Colab
#!pip install -q diffusers transformers accelerate omegaconf
```
## Load model checkpoints
Model weights may be stored in separate subfolders on the Hub or locally, in which case, you should use the [`~StableDiffusionXLPipeline.from_pretrained`] method:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/sdxl-turbo", torch_dtype=torch.float16, variant="fp16")
pipeline = pipeline.to("cuda")
```
You can also use the [`~StableDiffusionXLPipeline.from_single_file`] method to load a model checkpoint stored in a single file format (`.ckpt` or `.safetensors`) from the Hub or locally:
```py
from diffusers import StableDiffusionXLPipeline
import torch
pipeline = StableDiffusionXLPipeline.from_single_file(
"https://huggingface.co/stabilityai/sdxl-turbo/blob/main/sd_xl_turbo_1.0_fp16.safetensors", torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
```
## Text-to-image
For text-to-image, pass a text prompt. By default, SDXL Turbo generates a 512x512 image, and that resolution gives the best results. You can try setting the `height` and `width` parameters to 768x768 or 1024x1024, but you should expect quality degradations when doing so.
Make sure to set `guidance_scale` to 0.0 to disable, as the model was trained without it. A single inference step is enough to generate high quality images.
Increasing the number of steps to 2, 3 or 4 should improve image quality.
```py
from diffusers import AutoPipelineForText2Image
import torch
pipeline_text2image = AutoPipelineForText2Image.from_pretrained("stabilityai/sdxl-turbo", torch_dtype=torch.float16, variant="fp16")
pipeline_text2image = pipeline_text2image.to("cuda")
prompt = "A cinematic shot of a baby racoon wearing an intricate italian priest robe."
image = pipeline_text2image(prompt=prompt, guidance_scale=0.0, num_inference_steps=1).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/sdxl-turbo-text2img.png" alt="generated image of a racoon in a robe"/>
</div>
## Image-to-image
For image-to-image generation, make sure that `num_inference_steps * strength` is larger or equal to 1.
The image-to-image pipeline will run for `int(num_inference_steps * strength)` steps, e.g. `0.5 * 2.0 = 1` step in
our example below.
```py
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image, make_image_grid
# use from_pipe to avoid consuming additional memory when loading a checkpoint
pipeline = AutoPipelineForImage2Image.from_pipe(pipeline_text2image).to("cuda")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/cat.png")
init_image = init_image.resize((512, 512))
prompt = "cat wizard, gandalf, lord of the rings, detailed, fantasy, cute, adorable, Pixar, Disney, 8k"
image = pipeline(prompt, image=init_image, strength=0.5, guidance_scale=0.0, num_inference_steps=2).images[0]
make_image_grid([init_image, image], rows=1, cols=2)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/sdxl-turbo-img2img.png" alt="Image-to-image generation sample using SDXL Turbo"/>
</div>
## Speed-up SDXL Turbo even more
- Compile the UNet if you are using PyTorch version 2 or better. The first inference run will be very slow, but subsequent ones will be much faster.
```py
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
```
- When using the default VAE, keep it in `float32` to avoid costly `dtype` conversions before and after each generation. You only need to do this one before your first generation:
```py
pipe.upcast_vae()
```
As an alternative, you can also use a [16-bit VAE](https://huggingface.co/madebyollin/sdxl-vae-fp16-fix) created by community member [`@madebyollin`](https://huggingface.co/madebyollin) that does not need to be upcasted to `float32`.
+134
View File
@@ -0,0 +1,134 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Video Diffusion
[[open-in-colab]]
[Stable Video Diffusion](https://static1.squarespace.com/static/6213c340453c3f502425776e/t/655ce779b9d47d342a93c890/1700587395994/stable_video_diffusion.pdf) is a powerful image-to-video generation model that can generate high resolution (576x1024) 2-4 second videos conditioned on the input image.
This guide will show you how to use SVD to short generate videos from images.
Before you begin, make sure you have the following libraries installed:
```py
!pip install -q -U diffusers transformers accelerate
```
## Image to Video Generation
The are two variants of SVD. [SVD](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid)
and [SVD-XT](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid-xt). The svd checkpoint is trained to generate 14 frames and the svd-xt checkpoint is further
finetuned to generate 25 frames.
We will use the `svd-xt` checkpoint for this guide.
```python
import torch
from diffusers import StableVideoDiffusionPipeline
from diffusers.utils import load_image, export_to_video
pipe = StableVideoDiffusionPipeline.from_pretrained(
"stabilityai/stable-video-diffusion-img2vid-xt", torch_dtype=torch.float16, variant="fp16"
)
pipe.enable_model_cpu_offload()
# Load the conditioning image
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket.png?download=true")
image = image.resize((1024, 576))
generator = torch.manual_seed(42)
frames = pipe(image, decode_chunk_size=8, generator=generator).frames[0]
export_to_video(frames, "generated.mp4", fps=7)
```
<video controls width="1024" height="576">
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket_generated.webm" type="video/webm" />
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket_generated.mp4" type="video/mp4" />
</video>
<Tip>
Since generating videos is more memory intensive we can use the `decode_chunk_size` argument to control how many frames are decoded at once. This will reduce the memory usage. It's recommended to tweak this value based on your GPU memory.
Setting `decode_chunk_size=1` will decode one frame at a time and will use the least amount of memory but the video might have some flickering.
Additionally, we also use [model cpu offloading](../../optimization/memory#model-offloading) to reduce the memory usage.
</Tip>
### Torch.compile
You can achieve a 20-25% speed-up at the expense of slightly increased memory by compiling the UNet as follows:
```diff
- pipe.enable_model_cpu_offload()
+ pipe.to("cuda")
+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
```
### Low-memory
Video generation is very memory intensive as we have to essentially generate `num_frames` all at once. The mechanism is very comparable to text-to-image generation with a high batch size. To reduce the memory requirement you have multiple options. The following options trade inference speed against lower memory requirement:
- enable model offloading: Each component of the pipeline is offloaded to CPU once it's not needed anymore.
- enable feed-forward chunking: The feed-forward layer runs in a loop instead of running with a single huge feed-forward batch size
- reduce `decode_chunk_size`: This means that the VAE decodes frames in chunks instead of decoding them all together. **Note**: In addition to leading to a small slowdown, this method also slightly leads to video quality deterioration
You can enable them as follows:
```diff
-pipe.enable_model_cpu_offload()
-frames = pipe(image, decode_chunk_size=8, generator=generator).frames[0]
+pipe.enable_model_cpu_offload()
+pipe.unet.enable_forward_chunking()
+frames = pipe(image, decode_chunk_size=2, generator=generator, num_frames=25).frames[0]
```
Including all these tricks should lower the memory requirement to less than 8GB VRAM.
### Micro-conditioning
Along with conditioning image Stable Diffusion Video also allows providing micro-conditioning that allows more control over the generated video.
It accepts the following arguments:
- `fps`: The frames per second of the generated video.
- `motion_bucket_id`: The motion bucket id to use for the generated video. This can be used to control the motion of the generated video. Increasing the motion bucket id will increase the motion of the generated video.
- `noise_aug_strength`: The amount of noise added to the conditioning image. The higher the values the less the video will resemble the conditioning image. Increasing this value will also increase the motion of the generated video.
Here is an example of using micro-conditioning to generate a video with more motion.
```python
import torch
from diffusers import StableVideoDiffusionPipeline
from diffusers.utils import load_image, export_to_video
pipe = StableVideoDiffusionPipeline.from_pretrained(
"stabilityai/stable-video-diffusion-img2vid-xt", torch_dtype=torch.float16, variant="fp16"
)
pipe.enable_model_cpu_offload()
# Load the conditioning image
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket.png?download=true")
image = image.resize((1024, 576))
generator = torch.manual_seed(42)
frames = pipe(image, decode_chunk_size=8, generator=generator, motion_bucket_id=180, noise_aug_strength=0.1).frames[0]
export_to_video(frames, "generated.mp4", fps=7)
```
<video width="1024" height="576" controls>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket_generated_motion.mp4" type="video/mp4">
</video>
@@ -14,54 +14,41 @@ specific language governing permissions and limitations under the License.
[[open-in-colab]]
Unconditional image generation is a relatively straightforward task. The model only generates images - without any additional context like text or an image - resembling the training data it was trained on.
Unconditional image generation generates images that look like a random sample from the training data the model was trained on because the denoising process is not guided by any additional context like text or image.
The [`DiffusionPipeline`] is the easiest way to use a pre-trained diffusion system for inference.
To get started, use the [`DiffusionPipeline`] to load the [anton-l/ddpm-butterflies-128](https://huggingface.co/anton-l/ddpm-butterflies-128) checkpoint to generate images of butterflies. The [`DiffusionPipeline`] downloads and caches all the model components required to generate an image.
Start by creating an instance of [`DiffusionPipeline`] and specify which pipeline checkpoint you would like to download.
You can use any of the 🧨 Diffusers [checkpoints](https://huggingface.co/models?library=diffusers&sort=downloads) from the Hub (the checkpoint you'll use generates images of butterflies).
<Tip>
💡 Want to train your own unconditional image generation model? Take a look at the training [guide](../training/unconditional_training) to learn how to generate your own images.
</Tip>
In this guide, you'll use [`DiffusionPipeline`] for unconditional image generation with [DDPM](https://arxiv.org/abs/2006.11239):
```python
```py
from diffusers import DiffusionPipeline
generator = DiffusionPipeline.from_pretrained("anton-l/ddpm-butterflies-128", use_safetensors=True)
```
The [`DiffusionPipeline`] downloads and caches all modeling, tokenization, and scheduling components.
Because the model consists of roughly 1.4 billion parameters, we strongly recommend running it on a GPU.
You can move the generator object to a GPU, just like you would in PyTorch:
```python
generator.to("cuda")
```
Now you can use the `generator` to generate an image:
```python
generator = DiffusionPipeline.from_pretrained("anton-l/ddpm-butterflies-128").to("cuda")
image = generator().images[0]
image
```
The output is by default wrapped into a [`PIL.Image`](https://pillow.readthedocs.io/en/stable/reference/Image.html?highlight=image#the-image-class) object.
<Tip>
You can save the image by calling:
Want to generate images of something else? Take a look at the training [guide](../training/unconditional_training) to learn how to train a model to generate your own images.
```python
</Tip>
The output image is a [`PIL.Image`](https://pillow.readthedocs.io/en/stable/reference/Image.html?highlight=image#the-image-class) object that can be saved:
```py
image.save("generated_image.png")
```
Try out the Spaces below, and feel free to play around with the inference steps parameter to see how it affects the image quality!
You can also try experimenting with the `num_inference_steps` parameter, which controls the number of denoising steps. More denoising steps typically produce higher quality images, but it'll take longer to generate. Feel free to play around with this parameter to see how it affects the image quality.
```py
image = generator(num_inference_steps=100).images[0]
image
```
Try out the Space below to generate an image of a butterfly!
<iframe
src="https://stevhliu-ddpm-butterflies-128.hf.space"
src="https://stevhliu-unconditional-image-generation.hf.space"
frameborder="0"
width="850"
height="500"
+1
View File
@@ -96,3 +96,4 @@ specific language governing permissions and limitations under the License.
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
| [stable_diffusion_ldm3d](./api/pipelines/stable_diffusion/ldm3d_diffusion) | [LDM3D: Latent Diffusion Model for 3D](https://arxiv.org/abs/2305.10853) | Text to Image and Depth Generation |
| [stable_diffusion_upscaler_ldm3d](./api/pipelines/stable_diffusion/ldm3d_diffusion) | [LDM3D-VR: Latent Diffusion Model for 3D VR](https://arxiv.org/pdf/2311.03226) | Image and Depth Upscaling |
+4 -6
View File
@@ -18,8 +18,7 @@ limitations under the License.
Diffusers examples are a collection of scripts to demonstrate how to effectively use the `diffusers` library
for a variety of use cases involving training or fine-tuning.
**Note**: If you are looking for **official** examples on how to use `diffusers` for inference,
please have a look at [src/diffusers/pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines).
**Note**: If you are looking for **official** examples on how to use `diffusers` for inference, please have a look at [src/diffusers/pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines).
Our examples aspire to be **self-contained**, **easy-to-tweak**, **beginner-friendly** and for **one-purpose-only**.
More specifically, this means:
@@ -27,11 +26,10 @@ More specifically, this means:
- **Self-contained**: An example script shall only depend on "pip-install-able" Python packages that can be found in a `requirements.txt` file. Example scripts shall **not** depend on any local files. This means that one can simply download an example script, *e.g.* [train_unconditional.py](https://github.com/huggingface/diffusers/blob/main/examples/unconditional_image_generation/train_unconditional.py), install the required dependencies, *e.g.* [requirements.txt](https://github.com/huggingface/diffusers/blob/main/examples/unconditional_image_generation/requirements.txt) and execute the example script.
- **Easy-to-tweak**: While we strive to present as many use cases as possible, the example scripts are just that - examples. It is expected that they won't work out-of-the box on your specific problem and that you will be required to change a few lines of code to adapt them to your needs. To help you with that, most of the examples fully expose the preprocessing of the data and the training loop to allow you to tweak and edit them as required.
- **Beginner-friendly**: We do not aim for providing state-of-the-art training scripts for the newest models, but rather examples that can be used as a way to better understand diffusion models and how to use them with the `diffusers` library. We often purposefully leave out certain state-of-the-art methods if we consider them too complex for beginners.
- **One-purpose-only**: Examples should show one task and one task only. Even if a task is from a modeling
point of view very similar, *e.g.* image super-resolution and image modification tend to use the same model and training method, we want examples to showcase only one task to keep them as readable and easy-to-understand as possible.
- **One-purpose-only**: Examples should show one task and one task only. Even if a task is from a modeling point of view very similar, *e.g.* image super-resolution and image modification tend to use the same model and training method, we want examples to showcase only one task to keep them as readable and easy-to-understand as possible.
We provide **official** examples that cover the most popular tasks of diffusion models.
*Official* examples are **actively** maintained by the `diffusers` maintainers and we try to rigorously follow our example philosophy as defined above.
*Official* examples are **actively** maintained by the `diffusers` maintainers and we try to rigorously follow our example philosophy as defined above.
If you feel like another important example should exist, we are more than happy to welcome a [Feature Request](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feature_request.md&title=) or directly a [Pull Request](https://github.com/huggingface/diffusers/compare) from you!
Training examples show how to pretrain or fine-tune diffusion models for a variety of tasks. Currently we support:
@@ -39,7 +37,7 @@ Training examples show how to pretrain or fine-tune diffusion models for a varie
| Task | 🤗 Accelerate | 🤗 Datasets | Colab
|---|---|:---:|:---:|
| [**Unconditional Image Generation**](./unconditional_image_generation) | ✅ | ✅ | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [**Text-to-Image fine-tuning**](./text_to_image) | ✅ | ✅ |
| [**Text-to-Image fine-tuning**](./text_to_image) | ✅ | ✅ |
| [**Textual Inversion**](./textual_inversion) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_textual_inversion_training.ipynb)
| [**Dreambooth**](./dreambooth) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_dreambooth_training.ipynb)
| [**ControlNet**](./controlnet) | ✅ | ✅ | -
File diff suppressed because it is too large Load Diff
File diff suppressed because it is too large Load Diff
+8 -6
View File
@@ -5,10 +5,11 @@ from typing import Dict, List, Union
import safetensors.torch
import torch
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from diffusers import DiffusionPipeline, __version__
from diffusers.schedulers.scheduling_utils import SCHEDULER_CONFIG_NAME
from diffusers.utils import CONFIG_NAME, DIFFUSERS_CACHE, ONNX_WEIGHTS_NAME, WEIGHTS_NAME
from diffusers.utils import CONFIG_NAME, ONNX_WEIGHTS_NAME, WEIGHTS_NAME
class CheckpointMergerPipeline(DiffusionPipeline):
@@ -57,6 +58,7 @@ class CheckpointMergerPipeline(DiffusionPipeline):
return (temp_dict, meta_keys)
@torch.no_grad()
@validate_hf_hub_args
def merge(self, pretrained_model_name_or_path_list: List[Union[str, os.PathLike]], **kwargs):
"""
Returns a new pipeline object of the class 'DiffusionPipeline' with the merged checkpoints(weights) of the models passed
@@ -69,7 +71,7 @@ class CheckpointMergerPipeline(DiffusionPipeline):
**kwargs:
Supports all the default DiffusionPipeline.get_config_dict kwargs viz..
cache_dir, resume_download, force_download, proxies, local_files_only, use_auth_token, revision, torch_dtype, device_map.
cache_dir, resume_download, force_download, proxies, local_files_only, token, revision, torch_dtype, device_map.
alpha - The interpolation parameter. Ranges from 0 to 1. It affects the ratio in which the checkpoints are merged. A 0.8 alpha
would mean that the first model checkpoints would affect the final result far less than an alpha of 0.2
@@ -81,12 +83,12 @@ class CheckpointMergerPipeline(DiffusionPipeline):
"""
# Default kwargs from DiffusionPipeline
cache_dir = kwargs.pop("cache_dir", DIFFUSERS_CACHE)
cache_dir = kwargs.pop("cache_dir", None)
resume_download = kwargs.pop("resume_download", False)
force_download = kwargs.pop("force_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", False)
use_auth_token = kwargs.pop("use_auth_token", None)
token = kwargs.pop("token", None)
revision = kwargs.pop("revision", None)
torch_dtype = kwargs.pop("torch_dtype", None)
device_map = kwargs.pop("device_map", None)
@@ -123,7 +125,7 @@ class CheckpointMergerPipeline(DiffusionPipeline):
force_download=force_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
token=token,
revision=revision,
)
config_dicts.append(config_dict)
@@ -159,7 +161,7 @@ class CheckpointMergerPipeline(DiffusionPipeline):
resume_download=resume_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
token=token,
revision=revision,
allow_patterns=allow_patterns,
user_agent=user_agent,
+466
View File
@@ -0,0 +1,466 @@
# Copyright 2023 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
from math import pi
from typing import Callable, List, Optional, Tuple, Union
import numpy as np
import torch
from PIL import Image
from diffusers import DDPMScheduler, DiffusionPipeline, ImagePipelineOutput, UNet2DModel
from diffusers.utils.torch_utils import randn_tensor
class DPSPipeline(DiffusionPipeline):
r"""
Pipeline for Diffusion Posterior Sampling.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
Parameters:
unet ([`UNet2DModel`]):
A `UNet2DModel` to denoise the encoded image latents.
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image. Can be one of
[`DDPMScheduler`], or [`DDIMScheduler`].
"""
model_cpu_offload_seq = "unet"
def __init__(self, unet, scheduler):
super().__init__()
self.register_modules(unet=unet, scheduler=scheduler)
@torch.no_grad()
def __call__(
self,
measurement: torch.Tensor,
operator: torch.nn.Module,
loss_fn: Callable[[torch.Tensor, torch.Tensor], torch.Tensor],
batch_size: int = 1,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
num_inference_steps: int = 1000,
output_type: Optional[str] = "pil",
return_dict: bool = True,
zeta: float = 0.3,
) -> Union[ImagePipelineOutput, Tuple]:
r"""
The call function to the pipeline for generation.
Args:
measurement (`torch.Tensor`, *required*):
A 'torch.Tensor', the corrupted image
operator (`torch.nn.Module`, *required*):
A 'torch.nn.Module', the operator generating the corrupted image
loss_fn (`Callable[[torch.Tensor, torch.Tensor], torch.Tensor]`, *required*):
A 'Callable[[torch.Tensor, torch.Tensor], torch.Tensor]', the loss function used
between the measurements, for most of the cases using RMSE is fine.
batch_size (`int`, *optional*, defaults to 1):
The number of images to generate.
generator (`torch.Generator`, *optional*):
A [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make
generation deterministic.
num_inference_steps (`int`, *optional*, defaults to 1000):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generated image. Choose between `PIL.Image` or `np.array`.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~pipelines.ImagePipelineOutput`] instead of a plain tuple.
Example:
```py
>>> from diffusers import DDPMPipeline
>>> # load model and scheduler
>>> pipe = DDPMPipeline.from_pretrained("google/ddpm-cat-256")
>>> # run pipeline in inference (sample random noise and denoise)
>>> image = pipe().images[0]
>>> # save image
>>> image.save("ddpm_generated_image.png")
```
Returns:
[`~pipelines.ImagePipelineOutput`] or `tuple`:
If `return_dict` is `True`, [`~pipelines.ImagePipelineOutput`] is returned, otherwise a `tuple` is
returned where the first element is a list with the generated images
"""
# Sample gaussian noise to begin loop
if isinstance(self.unet.config.sample_size, int):
image_shape = (
batch_size,
self.unet.config.in_channels,
self.unet.config.sample_size,
self.unet.config.sample_size,
)
else:
image_shape = (batch_size, self.unet.config.in_channels, *self.unet.config.sample_size)
if self.device.type == "mps":
# randn does not work reproducibly on mps
image = randn_tensor(image_shape, generator=generator)
image = image.to(self.device)
else:
image = randn_tensor(image_shape, generator=generator, device=self.device)
# set step values
self.scheduler.set_timesteps(num_inference_steps)
for t in self.progress_bar(self.scheduler.timesteps):
with torch.enable_grad():
# 1. predict noise model_output
image = image.requires_grad_()
model_output = self.unet(image, t).sample
# 2. compute previous image x'_{t-1} and original prediction x0_{t}
scheduler_out = self.scheduler.step(model_output, t, image, generator=generator)
image_pred, origi_pred = scheduler_out.prev_sample, scheduler_out.pred_original_sample
# 3. compute y'_t = f(x0_{t})
measurement_pred = operator(origi_pred)
# 4. compute loss = d(y, y'_t-1)
loss = loss_fn(measurement, measurement_pred)
loss.backward()
print("distance: {0:.4f}".format(loss.item()))
with torch.no_grad():
image_pred = image_pred - zeta * image.grad
image = image_pred.detach()
image = (image / 2 + 0.5).clamp(0, 1)
image = image.cpu().permute(0, 2, 3, 1).numpy()
if output_type == "pil":
image = self.numpy_to_pil(image)
if not return_dict:
return (image,)
return ImagePipelineOutput(images=image)
if __name__ == "__main__":
import scipy
from torch import nn
from torchvision.utils import save_image
# defining the operators f(.) of y = f(x)
# super-resolution operator
class SuperResolutionOperator(nn.Module):
def __init__(self, in_shape, scale_factor):
super().__init__()
# Resizer local class, do not use outiside the SR operator class
class Resizer(nn.Module):
def __init__(self, in_shape, scale_factor=None, output_shape=None, kernel=None, antialiasing=True):
super(Resizer, self).__init__()
# First standardize values and fill missing arguments (if needed) by deriving scale from output shape or vice versa
scale_factor, output_shape = self.fix_scale_and_size(in_shape, output_shape, scale_factor)
# Choose interpolation method, each method has the matching kernel size
def cubic(x):
absx = np.abs(x)
absx2 = absx**2
absx3 = absx**3
return (1.5 * absx3 - 2.5 * absx2 + 1) * (absx <= 1) + (
-0.5 * absx3 + 2.5 * absx2 - 4 * absx + 2
) * ((1 < absx) & (absx <= 2))
def lanczos2(x):
return (
(np.sin(pi * x) * np.sin(pi * x / 2) + np.finfo(np.float32).eps)
/ ((pi**2 * x**2 / 2) + np.finfo(np.float32).eps)
) * (abs(x) < 2)
def box(x):
return ((-0.5 <= x) & (x < 0.5)) * 1.0
def lanczos3(x):
return (
(np.sin(pi * x) * np.sin(pi * x / 3) + np.finfo(np.float32).eps)
/ ((pi**2 * x**2 / 3) + np.finfo(np.float32).eps)
) * (abs(x) < 3)
def linear(x):
return (x + 1) * ((-1 <= x) & (x < 0)) + (1 - x) * ((0 <= x) & (x <= 1))
method, kernel_width = {
"cubic": (cubic, 4.0),
"lanczos2": (lanczos2, 4.0),
"lanczos3": (lanczos3, 6.0),
"box": (box, 1.0),
"linear": (linear, 2.0),
None: (cubic, 4.0), # set default interpolation method as cubic
}.get(kernel)
# Antialiasing is only used when downscaling
antialiasing *= np.any(np.array(scale_factor) < 1)
# Sort indices of dimensions according to scale of each dimension. since we are going dim by dim this is efficient
sorted_dims = np.argsort(np.array(scale_factor))
self.sorted_dims = [int(dim) for dim in sorted_dims if scale_factor[dim] != 1]
# Iterate over dimensions to calculate local weights for resizing and resize each time in one direction
field_of_view_list = []
weights_list = []
for dim in self.sorted_dims:
# for each coordinate (along 1 dim), calculate which coordinates in the input image affect its result and the
# weights that multiply the values there to get its result.
weights, field_of_view = self.contributions(
in_shape[dim], output_shape[dim], scale_factor[dim], method, kernel_width, antialiasing
)
# convert to torch tensor
weights = torch.tensor(weights.T, dtype=torch.float32)
# We add singleton dimensions to the weight matrix so we can multiply it with the big tensor we get for
# tmp_im[field_of_view.T], (bsxfun style)
weights_list.append(
nn.Parameter(
torch.reshape(weights, list(weights.shape) + (len(scale_factor) - 1) * [1]),
requires_grad=False,
)
)
field_of_view_list.append(
nn.Parameter(
torch.tensor(field_of_view.T.astype(np.int32), dtype=torch.long), requires_grad=False
)
)
self.field_of_view = nn.ParameterList(field_of_view_list)
self.weights = nn.ParameterList(weights_list)
def forward(self, in_tensor):
x = in_tensor
# Use the affecting position values and the set of weights to calculate the result of resizing along this 1 dim
for dim, fov, w in zip(self.sorted_dims, self.field_of_view, self.weights):
# To be able to act on each dim, we swap so that dim 0 is the wanted dim to resize
x = torch.transpose(x, dim, 0)
# This is a bit of a complicated multiplication: x[field_of_view.T] is a tensor of order image_dims+1.
# for each pixel in the output-image it matches the positions the influence it from the input image (along 1 dim
# only, this is why it only adds 1 dim to 5the shape). We then multiply, for each pixel, its set of positions with
# the matching set of weights. we do this by this big tensor element-wise multiplication (MATLAB bsxfun style:
# matching dims are multiplied element-wise while singletons mean that the matching dim is all multiplied by the
# same number
x = torch.sum(x[fov] * w, dim=0)
# Finally we swap back the axes to the original order
x = torch.transpose(x, dim, 0)
return x
def fix_scale_and_size(self, input_shape, output_shape, scale_factor):
# First fixing the scale-factor (if given) to be standardized the function expects (a list of scale factors in the
# same size as the number of input dimensions)
if scale_factor is not None:
# By default, if scale-factor is a scalar we assume 2d resizing and duplicate it.
if np.isscalar(scale_factor) and len(input_shape) > 1:
scale_factor = [scale_factor, scale_factor]
# We extend the size of scale-factor list to the size of the input by assigning 1 to all the unspecified scales
scale_factor = list(scale_factor)
scale_factor = [1] * (len(input_shape) - len(scale_factor)) + scale_factor
# Fixing output-shape (if given): extending it to the size of the input-shape, by assigning the original input-size
# to all the unspecified dimensions
if output_shape is not None:
output_shape = list(input_shape[len(output_shape) :]) + list(np.uint(np.array(output_shape)))
# Dealing with the case of non-give scale-factor, calculating according to output-shape. note that this is
# sub-optimal, because there can be different scales to the same output-shape.
if scale_factor is None:
scale_factor = 1.0 * np.array(output_shape) / np.array(input_shape)
# Dealing with missing output-shape. calculating according to scale-factor
if output_shape is None:
output_shape = np.uint(np.ceil(np.array(input_shape) * np.array(scale_factor)))
return scale_factor, output_shape
def contributions(self, in_length, out_length, scale, kernel, kernel_width, antialiasing):
# This function calculates a set of 'filters' and a set of field_of_view that will later on be applied
# such that each position from the field_of_view will be multiplied with a matching filter from the
# 'weights' based on the interpolation method and the distance of the sub-pixel location from the pixel centers
# around it. This is only done for one dimension of the image.
# When anti-aliasing is activated (default and only for downscaling) the receptive field is stretched to size of
# 1/sf. this means filtering is more 'low-pass filter'.
fixed_kernel = (lambda arg: scale * kernel(scale * arg)) if antialiasing else kernel
kernel_width *= 1.0 / scale if antialiasing else 1.0
# These are the coordinates of the output image
out_coordinates = np.arange(1, out_length + 1)
# since both scale-factor and output size can be provided simulatneously, perserving the center of the image requires shifting
# the output coordinates. the deviation is because out_length doesn't necesary equal in_length*scale.
# to keep the center we need to subtract half of this deivation so that we get equal margins for boths sides and center is preserved.
shifted_out_coordinates = out_coordinates - (out_length - in_length * scale) / 2
# These are the matching positions of the output-coordinates on the input image coordinates.
# Best explained by example: say we have 4 horizontal pixels for HR and we downscale by SF=2 and get 2 pixels:
# [1,2,3,4] -> [1,2]. Remember each pixel number is the middle of the pixel.
# The scaling is done between the distances and not pixel numbers (the right boundary of pixel 4 is transformed to
# the right boundary of pixel 2. pixel 1 in the small image matches the boundary between pixels 1 and 2 in the big
# one and not to pixel 2. This means the position is not just multiplication of the old pos by scale-factor).
# So if we measure distance from the left border, middle of pixel 1 is at distance d=0.5, border between 1 and 2 is
# at d=1, and so on (d = p - 0.5). we calculate (d_new = d_old / sf) which means:
# (p_new-0.5 = (p_old-0.5) / sf) -> p_new = p_old/sf + 0.5 * (1-1/sf)
match_coordinates = shifted_out_coordinates / scale + 0.5 * (1 - 1 / scale)
# This is the left boundary to start multiplying the filter from, it depends on the size of the filter
left_boundary = np.floor(match_coordinates - kernel_width / 2)
# Kernel width needs to be enlarged because when covering has sub-pixel borders, it must 'see' the pixel centers
# of the pixels it only covered a part from. So we add one pixel at each side to consider (weights can zeroize them)
expanded_kernel_width = np.ceil(kernel_width) + 2
# Determine a set of field_of_view for each each output position, these are the pixels in the input image
# that the pixel in the output image 'sees'. We get a matrix whos horizontal dim is the output pixels (big) and the
# vertical dim is the pixels it 'sees' (kernel_size + 2)
field_of_view = np.squeeze(
np.int16(np.expand_dims(left_boundary, axis=1) + np.arange(expanded_kernel_width) - 1)
)
# Assign weight to each pixel in the field of view. A matrix whos horizontal dim is the output pixels and the
# vertical dim is a list of weights matching to the pixel in the field of view (that are specified in
# 'field_of_view')
weights = fixed_kernel(1.0 * np.expand_dims(match_coordinates, axis=1) - field_of_view - 1)
# Normalize weights to sum up to 1. be careful from dividing by 0
sum_weights = np.sum(weights, axis=1)
sum_weights[sum_weights == 0] = 1.0
weights = 1.0 * weights / np.expand_dims(sum_weights, axis=1)
# We use this mirror structure as a trick for reflection padding at the boundaries
mirror = np.uint(np.concatenate((np.arange(in_length), np.arange(in_length - 1, -1, step=-1))))
field_of_view = mirror[np.mod(field_of_view, mirror.shape[0])]
# Get rid of weights and pixel positions that are of zero weight
non_zero_out_pixels = np.nonzero(np.any(weights, axis=0))
weights = np.squeeze(weights[:, non_zero_out_pixels])
field_of_view = np.squeeze(field_of_view[:, non_zero_out_pixels])
# Final products are the relative positions and the matching weights, both are output_size X fixed_kernel_size
return weights, field_of_view
self.down_sample = Resizer(in_shape, 1 / scale_factor)
for param in self.parameters():
param.requires_grad = False
def forward(self, data, **kwargs):
return self.down_sample(data)
# Gaussian blurring operator
class GaussialBlurOperator(nn.Module):
def __init__(self, kernel_size, intensity):
super().__init__()
class Blurkernel(nn.Module):
def __init__(self, blur_type="gaussian", kernel_size=31, std=3.0):
super().__init__()
self.blur_type = blur_type
self.kernel_size = kernel_size
self.std = std
self.seq = nn.Sequential(
nn.ReflectionPad2d(self.kernel_size // 2),
nn.Conv2d(3, 3, self.kernel_size, stride=1, padding=0, bias=False, groups=3),
)
self.weights_init()
def forward(self, x):
return self.seq(x)
def weights_init(self):
if self.blur_type == "gaussian":
n = np.zeros((self.kernel_size, self.kernel_size))
n[self.kernel_size // 2, self.kernel_size // 2] = 1
k = scipy.ndimage.gaussian_filter(n, sigma=self.std)
k = torch.from_numpy(k)
self.k = k
for name, f in self.named_parameters():
f.data.copy_(k)
def update_weights(self, k):
if not torch.is_tensor(k):
k = torch.from_numpy(k)
for name, f in self.named_parameters():
f.data.copy_(k)
def get_kernel(self):
return self.k
self.kernel_size = kernel_size
self.conv = Blurkernel(blur_type="gaussian", kernel_size=kernel_size, std=intensity)
self.kernel = self.conv.get_kernel()
self.conv.update_weights(self.kernel.type(torch.float32))
for param in self.parameters():
param.requires_grad = False
def forward(self, data, **kwargs):
return self.conv(data)
def transpose(self, data, **kwargs):
return data
def get_kernel(self):
return self.kernel.view(1, 1, self.kernel_size, self.kernel_size)
# assuming the forward process y = f(x) is polluted by Gaussian noise, use l2 norm
def RMSELoss(yhat, y):
return torch.sqrt(torch.sum((yhat - y) ** 2))
# set up source image
src = Image.open("sample.png")
# read image into [1,3,H,W]
src = torch.from_numpy(np.array(src, dtype=np.float32)).permute(2, 0, 1)[None]
# normalize image to [-1,1]
src = (src / 127.5) - 1.0
src = src.to("cuda")
# set up operator and measurement
# operator = SuperResolutionOperator(in_shape=src.shape, scale_factor=4).to("cuda")
operator = GaussialBlurOperator(kernel_size=61, intensity=3.0).to("cuda")
measurement = operator(src)
# set up scheduler
scheduler = DDPMScheduler.from_pretrained("google/ddpm-celebahq-256")
scheduler.set_timesteps(1000)
# set up model
model = UNet2DModel.from_pretrained("google/ddpm-celebahq-256").to("cuda")
save_image((src + 1.0) / 2.0, "dps_src.png")
save_image((measurement + 1.0) / 2.0, "dps_mea.png")
# finally, the pipeline
dpspipe = DPSPipeline(model, scheduler)
image = dpspipe(
measurement=measurement,
operator=operator,
loss_fn=RMSELoss,
zeta=1.0,
).images[0]
image.save("dps_generated_image.png")
+620 -22
View File
@@ -16,6 +16,7 @@
import ast
import gc
import inspect
import math
import warnings
from collections.abc import Iterable
@@ -23,16 +24,29 @@ from typing import Any, Callable, Dict, List, Optional, Union
import torch
import torch.nn.functional as F
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from packaging import version
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.configuration_utils import FrozenDict
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor
from diffusers.loaders import FromSingleFileMixin, IPAdapterMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.attention import Attention, GatedSelfAttentionDense
from diffusers.models.attention_processor import AttnProcessor2_0
from diffusers.pipelines.stable_diffusion import StableDiffusionPipeline
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines import DiffusionPipeline
from diffusers.pipelines.stable_diffusion.pipeline_output import StableDiffusionPipelineOutput
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
from diffusers.schedulers import KarrasDiffusionSchedulers
from diffusers.utils import logging, replace_example_docstring
from diffusers.utils import (
USE_PEFT_BACKEND,
deprecate,
logging,
replace_example_docstring,
scale_lora_layers,
unscale_lora_layers,
)
from diffusers.utils.torch_utils import randn_tensor
EXAMPLE_DOC_STRING = """
@@ -44,6 +58,7 @@ EXAMPLE_DOC_STRING = """
>>> pipe = DiffusionPipeline.from_pretrained(
... "longlian/lmd_plus",
... custom_pipeline="llm_grounded_diffusion",
... custom_revision="main",
... variant="fp16", torch_dtype=torch.float16
... )
>>> pipe.enable_model_cpu_offload()
@@ -96,7 +111,12 @@ logger = logging.get_logger(__name__) # pylint: disable=invalid-name
# All keys in Stable Diffusion models: [('down', 0, 0, 0), ('down', 0, 1, 0), ('down', 1, 0, 0), ('down', 1, 1, 0), ('down', 2, 0, 0), ('down', 2, 1, 0), ('mid', 0, 0, 0), ('up', 1, 0, 0), ('up', 1, 1, 0), ('up', 1, 2, 0), ('up', 2, 0, 0), ('up', 2, 1, 0), ('up', 2, 2, 0), ('up', 3, 0, 0), ('up', 3, 1, 0), ('up', 3, 2, 0)]
# Note that the first up block is `UpBlock2D` rather than `CrossAttnUpBlock2D` and does not have attention. The last index is always 0 in our case since we have one `BasicTransformerBlock` in each `Transformer2DModel`.
DEFAULT_GUIDANCE_ATTN_KEYS = [("mid", 0, 0, 0), ("up", 1, 0, 0), ("up", 1, 1, 0), ("up", 1, 2, 0)]
DEFAULT_GUIDANCE_ATTN_KEYS = [
("mid", 0, 0, 0),
("up", 1, 0, 0),
("up", 1, 1, 0),
("up", 1, 2, 0),
]
def convert_attn_keys(key):
@@ -126,7 +146,15 @@ def scale_proportion(obj_box, H, W):
# Adapted from the parent class `AttnProcessor2_0`
class AttnProcessorWithHook(AttnProcessor2_0):
def __init__(self, attn_processor_key, hidden_size, cross_attention_dim, hook=None, fast_attn=True, enabled=True):
def __init__(
self,
attn_processor_key,
hidden_size,
cross_attention_dim,
hook=None,
fast_attn=True,
enabled=True,
):
super().__init__()
self.attn_processor_key = attn_processor_key
self.hidden_size = hidden_size
@@ -165,15 +193,16 @@ class AttnProcessorWithHook(AttnProcessor2_0):
if attn.group_norm is not None:
hidden_states = attn.group_norm(hidden_states.transpose(1, 2)).transpose(1, 2)
query = attn.to_q(hidden_states, scale=scale)
args = () if USE_PEFT_BACKEND else (scale,)
query = attn.to_q(hidden_states, *args)
if encoder_hidden_states is None:
encoder_hidden_states = hidden_states
elif attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
key = attn.to_k(encoder_hidden_states, scale=scale)
value = attn.to_v(encoder_hidden_states, scale=scale)
key = attn.to_k(encoder_hidden_states, *args)
value = attn.to_v(encoder_hidden_states, *args)
inner_dim = key.shape[-1]
head_dim = inner_dim // attn.heads
@@ -186,7 +215,13 @@ class AttnProcessorWithHook(AttnProcessor2_0):
if self.hook is not None and self.enabled:
# Call the hook with query, key, value, and attention maps
self.hook(self.attn_processor_key, query_batch_dim, key_batch_dim, value_batch_dim, attention_probs)
self.hook(
self.attn_processor_key,
query_batch_dim,
key_batch_dim,
value_batch_dim,
attention_probs,
)
if self.fast_attn:
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
@@ -202,7 +237,12 @@ class AttnProcessorWithHook(AttnProcessor2_0):
# the output of sdp = (batch, num_heads, seq_len, head_dim)
# TODO: add support for attn.scale when we move to Torch 2.1
hidden_states = F.scaled_dot_product_attention(
query, key, value, attn_mask=attention_mask, dropout_p=0.0, is_causal=False
query,
key,
value,
attn_mask=attention_mask,
dropout_p=0.0,
is_causal=False,
)
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
hidden_states = hidden_states.to(query.dtype)
@@ -211,7 +251,7 @@ class AttnProcessorWithHook(AttnProcessor2_0):
hidden_states = attn.batch_to_head_dim(hidden_states)
# linear proj
hidden_states = attn.to_out[0](hidden_states, scale=scale)
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
@@ -226,7 +266,9 @@ class AttnProcessorWithHook(AttnProcessor2_0):
return hidden_states
class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
class LLMGroundedDiffusionPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, IPAdapterMixin, FromSingleFileMixin
):
r"""
Pipeline for layout-grounded text-to-image generation using LLM-grounded Diffusion (LMD+): https://arxiv.org/pdf/2305.13655.pdf.
@@ -257,6 +299,11 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
Whether a safety checker is needed for this pipeline.
"""
model_cpu_offload_seq = "text_encoder->unet->vae"
_optional_components = ["safety_checker", "feature_extractor", "image_encoder"]
_exclude_from_cpu_offload = ["safety_checker"]
_callback_tensor_inputs = ["latents", "prompt_embeds", "negative_prompt_embeds"]
objects_text = "Objects: "
bg_prompt_text = "Background prompt: "
bg_prompt_text_no_trailing_space = bg_prompt_text.rstrip()
@@ -272,12 +319,91 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
scheduler: KarrasDiffusionSchedulers,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPImageProcessor,
image_encoder: CLIPVisionModelWithProjection = None,
requires_safety_checker: bool = True,
):
super().__init__(
vae, text_encoder, tokenizer, unet, scheduler, safety_checker, feature_extractor, requires_safety_checker
)
# This is copied from StableDiffusionPipeline, with hook initizations for LMD+.
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} is outdated. `steps_offset`"
f" should be set to 1 instead of {scheduler.config.steps_offset}. Please make sure "
"to update the config accordingly as leaving `steps_offset` might led to incorrect results"
" in future versions. If you have downloaded this checkpoint from the Hugging Face Hub,"
" it would be very nice if you could open a Pull request for the `scheduler/scheduler_config.json`"
" file"
)
deprecate("steps_offset!=1", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["steps_offset"] = 1
scheduler._internal_dict = FrozenDict(new_config)
if hasattr(scheduler.config, "clip_sample") and scheduler.config.clip_sample is True:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} has not set the configuration `clip_sample`."
" `clip_sample` should be set to False in the configuration file. Please make sure to update the"
" config accordingly as not setting `clip_sample` in the config might lead to incorrect results in"
" future versions. If you have downloaded this checkpoint from the Hugging Face Hub, it would be very"
" nice if you could open a Pull request for the `scheduler/scheduler_config.json` file"
)
deprecate("clip_sample not set", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["clip_sample"] = False
scheduler._internal_dict = FrozenDict(new_config)
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
version.parse(unet.config._diffusers_version).base_version
) < version.parse("0.9.0.dev0")
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
deprecation_message = (
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
" the `unet/config.json` file"
)
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(unet.config)
new_config["sample_size"] = 64
unet._internal_dict = FrozenDict(new_config)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
)
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor)
self.register_to_config(requires_safety_checker=requires_safety_checker)
# Initialize the attention hooks for LLM-grounded Diffusion
self.register_attn_hooks(unet)
self._saved_attn = None
@@ -464,7 +590,14 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
return token_map
def get_phrase_indices(self, prompt, phrases, token_map=None, add_suffix_if_not_found=False, verbose=False):
def get_phrase_indices(
self,
prompt,
phrases,
token_map=None,
add_suffix_if_not_found=False,
verbose=False,
):
for obj in phrases:
# Suffix the prompt with object name for attention guidance if object is not in the prompt, using "|" to separate the prompt and the suffix
if obj not in prompt:
@@ -485,7 +618,14 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
phrase_token_map_str = " ".join(phrase_token_map)
if verbose:
logger.info("Full str:", token_map_str, "Substr:", phrase_token_map_str, "Phrase:", phrases)
logger.info(
"Full str:",
token_map_str,
"Substr:",
phrase_token_map_str,
"Phrase:",
phrases,
)
# Count the number of token before substr
# The substring comes with a trailing space that needs to be removed by minus one in the index.
@@ -552,7 +692,15 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
return loss
def compute_ca_loss(self, saved_attn, bboxes, phrase_indices, guidance_attn_keys, verbose=False, **kwargs):
def compute_ca_loss(
self,
saved_attn,
bboxes,
phrase_indices,
guidance_attn_keys,
verbose=False,
**kwargs,
):
"""
The `saved_attn` is supposed to be passed to `save_attn_to_dict` in `cross_attention_kwargs` prior to computing ths loss.
`AttnProcessor` will put attention maps into the `save_attn_to_dict`.
@@ -605,6 +753,7 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
latents: Optional[torch.FloatTensor] = None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
ip_adapter_image: Optional[PipelineImageInput] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
@@ -662,6 +811,7 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided, `negative_prompt_embeds` are generated from the `negative_prompt` input argument.
ip_adapter_image: (`PipelineImageInput`, *optional*): Optional image input to work with IP Adapters.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generated image. Choose between `PIL.Image` or `np.array`.
return_dict (`bool`, *optional*, defaults to `True`):
@@ -724,9 +874,10 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
phrase_indices = []
prompt_parsed = []
for prompt_item in prompt:
phrase_indices_parsed_item, prompt_parsed_item = self.get_phrase_indices(
prompt_item, add_suffix_if_not_found=True
)
(
phrase_indices_parsed_item,
prompt_parsed_item,
) = self.get_phrase_indices(prompt_item, add_suffix_if_not_found=True)
phrase_indices.append(phrase_indices_parsed_item)
prompt_parsed.append(prompt_parsed_item)
prompt = prompt_parsed
@@ -759,6 +910,11 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
if do_classifier_free_guidance:
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds])
if ip_adapter_image is not None:
image_embeds, negative_image_embeds = self.encode_image(ip_adapter_image, device, num_images_per_prompt)
if self.do_classifier_free_guidance:
image_embeds = torch.cat([negative_image_embeds, image_embeds])
# 4. Prepare timesteps
self.scheduler.set_timesteps(num_inference_steps, device=device)
timesteps = self.scheduler.timesteps
@@ -801,7 +957,10 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
if n_objs:
cond_boxes[:n_objs] = torch.tensor(boxes)
text_embeddings = torch.zeros(
max_objs, self.unet.config.cross_attention_dim, device=device, dtype=self.text_encoder.dtype
max_objs,
self.unet.config.cross_attention_dim,
device=device,
dtype=self.text_encoder.dtype,
)
if n_objs:
text_embeddings[:n_objs] = _text_embeddings
@@ -833,6 +992,9 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 6.1 Add image embeds for IP-Adapter
added_cond_kwargs = {"image_embeds": image_embeds} if ip_adapter_image is not None else None
loss_attn = torch.tensor(10000.0)
# 7. Denoising loop
@@ -869,6 +1031,7 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
t,
encoder_hidden_states=prompt_embeds,
cross_attention_kwargs=cross_attention_kwargs,
added_cond_kwargs=added_cond_kwargs,
).sample
# perform guidance
@@ -1013,3 +1176,438 @@ class LLMGroundedDiffusionPipeline(StableDiffusionPipeline):
self.enable_attn_hook(enabled=False)
return latents, loss
# Below are methods copied from StableDiffusionPipeline
# The design choice of not inheriting from StableDiffusionPipeline is discussed here: https://github.com/huggingface/diffusers/pull/5993#issuecomment-1834258517
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_slicing
def enable_vae_slicing(self):
r"""
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
"""
self.vae.enable_slicing()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_slicing
def disable_vae_slicing(self):
r"""
Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_slicing()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_tiling
def enable_vae_tiling(self):
r"""
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow
processing larger images.
"""
self.vae.enable_tiling()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling
def disable_vae_tiling(self):
r"""
Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to
computing decoding in one step.
"""
self.vae.disable_tiling()
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._encode_prompt
def _encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
**kwargs,
):
deprecation_message = "`_encode_prompt()` is deprecated and it will be removed in a future version. Use `encode_prompt()` instead. Also, be aware that the output format changed from a concatenated tensor to a tuple."
deprecate("_encode_prompt()", "1.0.0", deprecation_message, standard_warn=False)
prompt_embeds_tuple = self.encode_prompt(
prompt=prompt,
device=device,
num_images_per_prompt=num_images_per_prompt,
do_classifier_free_guidance=do_classifier_free_guidance,
negative_prompt=negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
lora_scale=lora_scale,
**kwargs,
)
# concatenate for backwards comp
prompt_embeds = torch.cat([prompt_embeds_tuple[1], prompt_embeds_tuple[0]])
return prompt_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.encode_prompt
def encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
clip_skip: Optional[int] = None,
):
r"""
Encodes the prompt into text encoder hidden states.
Args:
prompt (`str` or `List[str]`, *optional*):
prompt to be encoded
device: (`torch.device`):
torch device
num_images_per_prompt (`int`):
number of images that should be generated per prompt
do_classifier_free_guidance (`bool`):
whether to use classifier free guidance or not
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. If not defined, one has to pass
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is
less than `1`).
prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not
provided, text embeddings will be generated from `prompt` input argument.
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input
argument.
lora_scale (`float`, *optional*):
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(self.text_encoder, lora_scale)
else:
scale_lora_layers(self.text_encoder, lora_scale)
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
if prompt_embeds is None:
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
prompt = self.maybe_convert_prompt(prompt, self.tokenizer)
text_inputs = self.tokenizer(
prompt,
padding="max_length",
max_length=self.tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
untruncated_ids = self.tokenizer(prompt, padding="longest", return_tensors="pt").input_ids
if untruncated_ids.shape[-1] >= text_input_ids.shape[-1] and not torch.equal(
text_input_ids, untruncated_ids
):
removed_text = self.tokenizer.batch_decode(
untruncated_ids[:, self.tokenizer.model_max_length - 1 : -1]
)
logger.warning(
"The following part of your input was truncated because CLIP can only handle sequences up to"
f" {self.tokenizer.model_max_length} tokens: {removed_text}"
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = text_inputs.attention_mask.to(device)
else:
attention_mask = None
if clip_skip is None:
prompt_embeds = self.text_encoder(text_input_ids.to(device), attention_mask=attention_mask)
prompt_embeds = prompt_embeds[0]
else:
prompt_embeds = self.text_encoder(
text_input_ids.to(device),
attention_mask=attention_mask,
output_hidden_states=True,
)
# Access the `hidden_states` first, that contains a tuple of
# all the hidden states from the encoder layers. Then index into
# the tuple to access the hidden states from the desired layer.
prompt_embeds = prompt_embeds[-1][-(clip_skip + 1)]
# We also need to apply the final LayerNorm here to not mess with the
# representations. The `last_hidden_states` that we typically use for
# obtaining the final prompt representations passes through the LayerNorm
# layer.
prompt_embeds = self.text_encoder.text_model.final_layer_norm(prompt_embeds)
if self.text_encoder is not None:
prompt_embeds_dtype = self.text_encoder.dtype
elif self.unet is not None:
prompt_embeds_dtype = self.unet.dtype
else:
prompt_embeds_dtype = prompt_embeds.dtype
prompt_embeds = prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
bs_embed, seq_len, _ = prompt_embeds.shape
# duplicate text embeddings for each generation per prompt, using mps friendly method
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1)
# get unconditional embeddings for classifier free guidance
if do_classifier_free_guidance and negative_prompt_embeds is None:
uncond_tokens: List[str]
if negative_prompt is None:
uncond_tokens = [""] * batch_size
elif prompt is not None and type(prompt) is not type(negative_prompt):
raise TypeError(
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !="
f" {type(prompt)}."
)
elif isinstance(negative_prompt, str):
uncond_tokens = [negative_prompt]
elif batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
else:
uncond_tokens = negative_prompt
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
uncond_tokens = self.maybe_convert_prompt(uncond_tokens, self.tokenizer)
max_length = prompt_embeds.shape[1]
uncond_input = self.tokenizer(
uncond_tokens,
padding="max_length",
max_length=max_length,
truncation=True,
return_tensors="pt",
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = uncond_input.attention_mask.to(device)
else:
attention_mask = None
negative_prompt_embeds = self.text_encoder(
uncond_input.input_ids.to(device),
attention_mask=attention_mask,
)
negative_prompt_embeds = negative_prompt_embeds[0]
if do_classifier_free_guidance:
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
seq_len = negative_prompt_embeds.shape[1]
negative_prompt_embeds = negative_prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
return prompt_embeds, negative_prompt_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.encode_image
def encode_image(self, image, device, num_images_per_prompt):
dtype = next(self.image_encoder.parameters()).dtype
if not isinstance(image, torch.Tensor):
image = self.feature_extractor(image, return_tensors="pt").pixel_values
image = image.to(device=device, dtype=dtype)
image_embeds = self.image_encoder(image).image_embeds
image_embeds = image_embeds.repeat_interleave(num_images_per_prompt, dim=0)
uncond_image_embeds = torch.zeros_like(image_embeds)
return image_embeds, uncond_image_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.run_safety_checker
def run_safety_checker(self, image, device, dtype):
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
safety_checker_input = self.feature_extractor(feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.decode_latents
def decode_latents(self, latents):
deprecation_message = "The decode_latents method is deprecated and will be removed in 1.0.0. Please use VaeImageProcessor.postprocess(...) instead"
deprecate("decode_latents", "1.0.0", deprecation_message, standard_warn=False)
latents = 1 / self.vae.config.scaling_factor * latents
image = self.vae.decode(latents, return_dict=False)[0]
image = (image / 2 + 0.5).clamp(0, 1)
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloat16
image = image.cpu().permute(0, 2, 3, 1).float().numpy()
return image
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs
def prepare_extra_step_kwargs(self, generator, eta):
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# check if the scheduler accepts generator
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
if accepts_generator:
extra_step_kwargs["generator"] = generator
return extra_step_kwargs
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
def prepare_latents(
self,
batch_size,
num_channels_latents,
height,
width,
dtype,
device,
generator,
latents=None,
):
shape = (
batch_size,
num_channels_latents,
height // self.vae_scale_factor,
width // self.vae_scale_factor,
)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
if latents is None:
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
else:
latents = latents.to(device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
return latents
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_freeu
def enable_freeu(self, s1: float, s2: float, b1: float, b2: float):
r"""Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497.
The suffixes after the scaling factors represent the stages where they are being applied.
Please refer to the [official repository](https://github.com/ChenyangSi/FreeU) for combinations of the values
that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL.
Args:
s1 (`float`):
Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to
mitigate "oversmoothing effect" in the enhanced denoising process.
s2 (`float`):
Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to
mitigate "oversmoothing effect" in the enhanced denoising process.
b1 (`float`): Scaling factor for stage 1 to amplify the contributions of backbone features.
b2 (`float`): Scaling factor for stage 2 to amplify the contributions of backbone features.
"""
if not hasattr(self, "unet"):
raise ValueError("The pipeline must have `unet` for using FreeU.")
self.unet.enable_freeu(s1=s1, s2=s2, b1=b1, b2=b2)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling
def disable_freeu(self):
"""Disables the FreeU mechanism if enabled."""
self.unet.disable_freeu()
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32):
"""
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
Args:
timesteps (`torch.Tensor`):
generate embedding vectors at these timesteps
embedding_dim (`int`, *optional*, defaults to 512):
dimension of the embeddings to generate
dtype:
data type of the generated embeddings
Returns:
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)`
"""
assert len(w.shape) == 1
w = w * 1000.0
half_dim = embedding_dim // 2
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1)
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb)
emb = w.to(dtype)[:, None] * emb[None, :]
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1)
if embedding_dim % 2 == 1: # zero pad
emb = torch.nn.functional.pad(emb, (0, 1))
assert emb.shape == (w.shape[0], embedding_dim)
return emb
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.guidance_scale
@property
def guidance_scale(self):
return self._guidance_scale
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.guidance_rescale
@property
def guidance_rescale(self):
return self._guidance_rescale
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.clip_skip
@property
def clip_skip(self):
return self._clip_skip
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.do_classifier_free_guidance
@property
def do_classifier_free_guidance(self):
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.cross_attention_kwargs
@property
def cross_attention_kwargs(self):
return self._cross_attention_kwargs
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.num_timesteps
@property
def num_timesteps(self):
return self._num_timesteps
+14 -10
View File
@@ -250,6 +250,7 @@ def get_weighted_text_embeddings_sdxl(
neg_prompt: str = "",
neg_prompt_2: str = None,
num_images_per_prompt: int = 1,
device: Optional[torch.device] = None,
):
"""
This function can process long prompt with weights, no length limitation
@@ -262,10 +263,13 @@ def get_weighted_text_embeddings_sdxl(
neg_prompt (str)
neg_prompt_2 (str)
num_images_per_prompt (int)
device (torch.device)
Returns:
prompt_embeds (torch.Tensor)
neg_prompt_embeds (torch.Tensor)
"""
device = device or pipe._execution_device
if prompt_2:
prompt = f"{prompt} {prompt_2}"
@@ -330,17 +334,17 @@ def get_weighted_text_embeddings_sdxl(
# get prompt embeddings one by one is not working.
for i in range(len(prompt_token_groups)):
# get positive prompt embeddings with weights
token_tensor = torch.tensor([prompt_token_groups[i]], dtype=torch.long, device=pipe.device)
weight_tensor = torch.tensor(prompt_weight_groups[i], dtype=torch.float16, device=pipe.device)
token_tensor = torch.tensor([prompt_token_groups[i]], dtype=torch.long, device=device)
weight_tensor = torch.tensor(prompt_weight_groups[i], dtype=torch.float16, device=device)
token_tensor_2 = torch.tensor([prompt_token_groups_2[i]], dtype=torch.long, device=pipe.device)
token_tensor_2 = torch.tensor([prompt_token_groups_2[i]], dtype=torch.long, device=device)
# use first text encoder
prompt_embeds_1 = pipe.text_encoder(token_tensor.to(pipe.device), output_hidden_states=True)
prompt_embeds_1 = pipe.text_encoder(token_tensor.to(device), output_hidden_states=True)
prompt_embeds_1_hidden_states = prompt_embeds_1.hidden_states[-2]
# use second text encoder
prompt_embeds_2 = pipe.text_encoder_2(token_tensor_2.to(pipe.device), output_hidden_states=True)
prompt_embeds_2 = pipe.text_encoder_2(token_tensor_2.to(device), output_hidden_states=True)
prompt_embeds_2_hidden_states = prompt_embeds_2.hidden_states[-2]
pooled_prompt_embeds = prompt_embeds_2[0]
@@ -357,16 +361,16 @@ def get_weighted_text_embeddings_sdxl(
embeds.append(token_embedding)
# get negative prompt embeddings with weights
neg_token_tensor = torch.tensor([neg_prompt_token_groups[i]], dtype=torch.long, device=pipe.device)
neg_token_tensor_2 = torch.tensor([neg_prompt_token_groups_2[i]], dtype=torch.long, device=pipe.device)
neg_weight_tensor = torch.tensor(neg_prompt_weight_groups[i], dtype=torch.float16, device=pipe.device)
neg_token_tensor = torch.tensor([neg_prompt_token_groups[i]], dtype=torch.long, device=device)
neg_token_tensor_2 = torch.tensor([neg_prompt_token_groups_2[i]], dtype=torch.long, device=device)
neg_weight_tensor = torch.tensor(neg_prompt_weight_groups[i], dtype=torch.float16, device=device)
# use first text encoder
neg_prompt_embeds_1 = pipe.text_encoder(neg_token_tensor.to(pipe.device), output_hidden_states=True)
neg_prompt_embeds_1 = pipe.text_encoder(neg_token_tensor.to(device), output_hidden_states=True)
neg_prompt_embeds_1_hidden_states = neg_prompt_embeds_1.hidden_states[-2]
# use second text encoder
neg_prompt_embeds_2 = pipe.text_encoder_2(neg_token_tensor_2.to(pipe.device), output_hidden_states=True)
neg_prompt_embeds_2 = pipe.text_encoder_2(neg_token_tensor_2.to(device), output_hidden_states=True)
neg_prompt_embeds_2_hidden_states = neg_prompt_embeds_2.hidden_states[-2]
negative_pooled_prompt_embeds = neg_prompt_embeds_2[0]
File diff suppressed because it is too large Load Diff
File diff suppressed because it is too large Load Diff
@@ -0,0 +1,772 @@
# Copyright 2023 The Intel Labs Team Authors and the HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import inspect
from typing import Any, Callable, Dict, List, Optional, Union
import numpy as np
import PIL
import torch
from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
from diffusers import DiffusionPipeline
from diffusers.image_processor import PipelineDepthInput, PipelineImageInput, VaeImageProcessorLDM3D
from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.models.lora import adjust_lora_scale_text_encoder
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d import LDM3DPipelineOutput
from diffusers.schedulers import DDPMScheduler, KarrasDiffusionSchedulers
from diffusers.utils import (
USE_PEFT_BACKEND,
deprecate,
logging,
scale_lora_layers,
unscale_lora_layers,
)
from diffusers.utils.torch_utils import randn_tensor
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
EXAMPLE_DOC_STRING = """
Examples:
```python
>>> from diffusers import StableDiffusionUpscaleLDM3DPipeline
>>> from PIL import Image
>>> from io import BytesIO
>>> import requests
>>> pipe = StableDiffusionUpscaleLDM3DPipeline.from_pretrained("Intel/ldm3d-sr")
>>> pipe = pipe.to("cuda")
>>> rgb_path = "https://huggingface.co/Intel/ldm3d-sr/resolve/main/lemons_ldm3d_rgb.jpg"
>>> depth_path = "https://huggingface.co/Intel/ldm3d-sr/resolve/main/lemons_ldm3d_depth.png"
>>> low_res_rgb = Image.open(BytesIO(requests.get(rgb_path).content)).convert("RGB")
>>> low_res_depth = Image.open(BytesIO(requests.get(depth_path).content)).convert("L")
>>> output = pipe(
... prompt="high quality high resolution uhd 4k image",
... rgb=low_res_rgb,
... depth=low_res_depth,
... num_inference_steps=50,
... target_res=[1024, 1024],
... )
>>> rgb_image, depth_image = output.rgb, output.depth
>>> rgb_image[0].save("hr_ldm3d_rgb.jpg")
>>> depth_image[0].save("hr_ldm3d_depth.png")
```
"""
class StableDiffusionUpscaleLDM3DPipeline(
DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
):
r"""
Pipeline for text-to-image and 3D generation using LDM3D.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods:
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings
- [`~loaders.LoraLoaderMixin.load_lora_weights`] for loading LoRA weights
- [`~loaders.LoraLoaderMixin.save_lora_weights`] for saving LoRA weights
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
text_encoder ([`~transformers.CLIPTextModel`]):
Frozen text-encoder ([clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14)).
tokenizer ([`~transformers.CLIPTokenizer`]):
A `CLIPTokenizer` to tokenize text.
unet ([`UNet2DConditionModel`]):
A `UNet2DConditionModel` to denoise the encoded image latents.
low_res_scheduler ([`SchedulerMixin`]):
A scheduler used to add initial noise to the low resolution conditioning image. It must be an instance of
[`DDPMScheduler`].
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor"]
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
low_res_scheduler: DDPMScheduler,
scheduler: KarrasDiffusionSchedulers,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPImageProcessor,
requires_safety_checker: bool = True,
watermarker: Optional[Any] = None,
max_noise_level: int = 350,
):
super().__init__()
if safety_checker is None and requires_safety_checker:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
if safety_checker is not None and feature_extractor is None:
raise ValueError(
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
low_res_scheduler=low_res_scheduler,
scheduler=scheduler,
safety_checker=safety_checker,
watermarker=watermarker,
feature_extractor=feature_extractor,
)
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
self.image_processor = VaeImageProcessorLDM3D(vae_scale_factor=self.vae_scale_factor, resample="bilinear")
# self.register_to_config(requires_safety_checker=requires_safety_checker)
self.register_to_config(max_noise_level=max_noise_level)
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d.StableDiffusionLDM3DPipeline._encode_prompt
def _encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
**kwargs,
):
deprecation_message = "`_encode_prompt()` is deprecated and it will be removed in a future version. Use `encode_prompt()` instead. Also, be aware that the output format changed from a concatenated tensor to a tuple."
deprecate("_encode_prompt()", "1.0.0", deprecation_message, standard_warn=False)
prompt_embeds_tuple = self.encode_prompt(
prompt=prompt,
device=device,
num_images_per_prompt=num_images_per_prompt,
do_classifier_free_guidance=do_classifier_free_guidance,
negative_prompt=negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
lora_scale=lora_scale,
**kwargs,
)
# concatenate for backwards comp
prompt_embeds = torch.cat([prompt_embeds_tuple[1], prompt_embeds_tuple[0]])
return prompt_embeds
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d.StableDiffusionLDM3DPipeline.encode_prompt
def encode_prompt(
self,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt=None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
lora_scale: Optional[float] = None,
clip_skip: Optional[int] = None,
):
r"""
Encodes the prompt into text encoder hidden states.
Args:
prompt (`str` or `List[str]`, *optional*):
prompt to be encoded
device: (`torch.device`):
torch device
num_images_per_prompt (`int`):
number of images that should be generated per prompt
do_classifier_free_guidance (`bool`):
whether to use classifier free guidance or not
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. If not defined, one has to pass
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is
less than `1`).
prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not
provided, text embeddings will be generated from `prompt` input argument.
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input
argument.
lora_scale (`float`, *optional*):
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
clip_skip (`int`, *optional*):
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that
the output of the pre-final layer will be used for computing the prompt embeddings.
"""
# set lora scale so that monkey patched LoRA
# function of text encoder can correctly access it
if lora_scale is not None and isinstance(self, LoraLoaderMixin):
self._lora_scale = lora_scale
# dynamically adjust the LoRA scale
if not USE_PEFT_BACKEND:
adjust_lora_scale_text_encoder(self.text_encoder, lora_scale)
else:
scale_lora_layers(self.text_encoder, lora_scale)
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
if prompt_embeds is None:
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
prompt = self.maybe_convert_prompt(prompt, self.tokenizer)
text_inputs = self.tokenizer(
prompt,
padding="max_length",
max_length=self.tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
untruncated_ids = self.tokenizer(prompt, padding="longest", return_tensors="pt").input_ids
if untruncated_ids.shape[-1] >= text_input_ids.shape[-1] and not torch.equal(
text_input_ids, untruncated_ids
):
removed_text = self.tokenizer.batch_decode(
untruncated_ids[:, self.tokenizer.model_max_length - 1 : -1]
)
logger.warning(
"The following part of your input was truncated because CLIP can only handle sequences up to"
f" {self.tokenizer.model_max_length} tokens: {removed_text}"
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = text_inputs.attention_mask.to(device)
else:
attention_mask = None
if clip_skip is None:
prompt_embeds = self.text_encoder(text_input_ids.to(device), attention_mask=attention_mask)
prompt_embeds = prompt_embeds[0]
else:
prompt_embeds = self.text_encoder(
text_input_ids.to(device), attention_mask=attention_mask, output_hidden_states=True
)
# Access the `hidden_states` first, that contains a tuple of
# all the hidden states from the encoder layers. Then index into
# the tuple to access the hidden states from the desired layer.
prompt_embeds = prompt_embeds[-1][-(clip_skip + 1)]
# We also need to apply the final LayerNorm here to not mess with the
# representations. The `last_hidden_states` that we typically use for
# obtaining the final prompt representations passes through the LayerNorm
# layer.
prompt_embeds = self.text_encoder.text_model.final_layer_norm(prompt_embeds)
if self.text_encoder is not None:
prompt_embeds_dtype = self.text_encoder.dtype
elif self.unet is not None:
prompt_embeds_dtype = self.unet.dtype
else:
prompt_embeds_dtype = prompt_embeds.dtype
prompt_embeds = prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
bs_embed, seq_len, _ = prompt_embeds.shape
# duplicate text embeddings for each generation per prompt, using mps friendly method
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1)
prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1)
# get unconditional embeddings for classifier free guidance
if do_classifier_free_guidance and negative_prompt_embeds is None:
uncond_tokens: List[str]
if negative_prompt is None:
uncond_tokens = [""] * batch_size
elif prompt is not None and type(prompt) is not type(negative_prompt):
raise TypeError(
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !="
f" {type(prompt)}."
)
elif isinstance(negative_prompt, str):
uncond_tokens = [negative_prompt]
elif batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
else:
uncond_tokens = negative_prompt
# textual inversion: procecss multi-vector tokens if necessary
if isinstance(self, TextualInversionLoaderMixin):
uncond_tokens = self.maybe_convert_prompt(uncond_tokens, self.tokenizer)
max_length = prompt_embeds.shape[1]
uncond_input = self.tokenizer(
uncond_tokens,
padding="max_length",
max_length=max_length,
truncation=True,
return_tensors="pt",
)
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
attention_mask = uncond_input.attention_mask.to(device)
else:
attention_mask = None
negative_prompt_embeds = self.text_encoder(
uncond_input.input_ids.to(device),
attention_mask=attention_mask,
)
negative_prompt_embeds = negative_prompt_embeds[0]
if do_classifier_free_guidance:
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
seq_len = negative_prompt_embeds.shape[1]
negative_prompt_embeds = negative_prompt_embeds.to(dtype=prompt_embeds_dtype, device=device)
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1)
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1)
if isinstance(self, LoraLoaderMixin) and USE_PEFT_BACKEND:
# Retrieve the original scale by scaling back the LoRA layers
unscale_lora_layers(self.text_encoder, lora_scale)
return prompt_embeds, negative_prompt_embeds
def run_safety_checker(self, image, device, dtype):
if self.safety_checker is None:
has_nsfw_concept = None
else:
if torch.is_tensor(image):
feature_extractor_input = self.image_processor.postprocess(image, output_type="pil")
else:
feature_extractor_input = self.image_processor.numpy_to_pil(image)
rgb_feature_extractor_input = feature_extractor_input[0]
safety_checker_input = self.feature_extractor(rgb_feature_extractor_input, return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
return image, has_nsfw_concept
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs
def prepare_extra_step_kwargs(self, generator, eta):
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# check if the scheduler accepts generator
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
if accepts_generator:
extra_step_kwargs["generator"] = generator
return extra_step_kwargs
def check_inputs(
self,
prompt,
image,
noise_level,
callback_steps,
negative_prompt=None,
prompt_embeds=None,
negative_prompt_embeds=None,
target_res=None,
):
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
if prompt is not None and prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to"
" only forward one of the two."
)
elif prompt is None and prompt_embeds is None:
raise ValueError(
"Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined."
)
elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if negative_prompt is not None and negative_prompt_embeds is not None:
raise ValueError(
f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:"
f" {negative_prompt_embeds}. Please make sure to only forward one of the two."
)
if prompt_embeds is not None and negative_prompt_embeds is not None:
if prompt_embeds.shape != negative_prompt_embeds.shape:
raise ValueError(
"`prompt_embeds` and `negative_prompt_embeds` must have the same shape when passed directly, but"
f" got: `prompt_embeds` {prompt_embeds.shape} != `negative_prompt_embeds`"
f" {negative_prompt_embeds.shape}."
)
if (
not isinstance(image, torch.Tensor)
and not isinstance(image, PIL.Image.Image)
and not isinstance(image, np.ndarray)
and not isinstance(image, list)
):
raise ValueError(
f"`image` has to be of type `torch.Tensor`, `np.ndarray`, `PIL.Image.Image` or `list` but is {type(image)}"
)
# verify batch size of prompt and image are same if image is a list or tensor or numpy array
if isinstance(image, list) or isinstance(image, torch.Tensor) or isinstance(image, np.ndarray):
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
if isinstance(image, list):
image_batch_size = len(image)
else:
image_batch_size = image.shape[0]
if batch_size != image_batch_size:
raise ValueError(
f"`prompt` has batch size {batch_size} and `image` has batch size {image_batch_size}."
" Please make sure that passed `prompt` matches the batch size of `image`."
)
# check noise level
if noise_level > self.config.max_noise_level:
raise ValueError(f"`noise_level` has to be <= {self.config.max_noise_level} but is {noise_level}")
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
shape = (batch_size, num_channels_latents, height, width)
if latents is None:
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
else:
if latents.shape != shape:
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {shape}")
latents = latents.to(device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
return latents
# def upcast_vae(self):
# dtype = self.vae.dtype
# self.vae.to(dtype=torch.float32)
# use_torch_2_0_or_xformers = isinstance(
# self.vae.decoder.mid_block.attentions[0].processor,
# (
# AttnProcessor2_0,
# XFormersAttnProcessor,
# LoRAXFormersAttnProcessor,
# LoRAAttnProcessor2_0,
# ),
# )
# # if xformers or torch_2_0 is used attention block does not need
# # to be in float32 which can save lots of memory
# if use_torch_2_0_or_xformers:
# self.vae.post_quant_conv.to(dtype)
# self.vae.decoder.conv_in.to(dtype)
# self.vae.decoder.mid_block.to(dtype)
@torch.no_grad()
def __call__(
self,
prompt: Union[str, List[str]] = None,
rgb: PipelineImageInput = None,
depth: PipelineDepthInput = None,
num_inference_steps: int = 75,
guidance_scale: float = 9.0,
noise_level: int = 20,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
latents: Optional[torch.FloatTensor] = None,
prompt_embeds: Optional[torch.FloatTensor] = None,
negative_prompt_embeds: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: int = 1,
cross_attention_kwargs: Optional[Dict[str, Any]] = None,
target_res: Optional[List[int]] = [1024, 1024],
):
r"""
The call function to the pipeline for generation.
Args:
prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide image generation. If not defined, you need to pass `prompt_embeds`.
image (`torch.FloatTensor`, `PIL.Image.Image`, `np.ndarray`, `List[torch.FloatTensor]`, `List[PIL.Image.Image]`, or `List[np.ndarray]`):
`Image` or tensor representing an image batch to be upscaled.
num_inference_steps (`int`, *optional*, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
guidance_scale (`float`, *optional*, defaults to 5.0):
A higher guidance scale value encourages the model to generate images closely linked to the text
`prompt` at the expense of lower image quality. Guidance scale is enabled when `guidance_scale > 1`.
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass `negative_prompt_embeds` instead. Ignored when not using guidance (`guidance_scale < 1`).
num_images_per_prompt (`int`, *optional*, defaults to 1):
The number of images to generate per prompt.
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) from the [DDIM](https://arxiv.org/abs/2010.02502) paper. Only applies
to the [`~schedulers.DDIMScheduler`], and is ignored in other schedulers.
generator (`torch.Generator` or `List[torch.Generator]`, *optional*):
A [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make
generation deterministic.
latents (`torch.FloatTensor`, *optional*):
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random `generator`.
prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the `prompt` input argument.
negative_prompt_embeds (`torch.FloatTensor`, *optional*):
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided, `negative_prompt_embeds` are generated from the `negative_prompt` input argument.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generated image. Choose between `PIL.Image` or `np.array`.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a
plain tuple.
callback (`Callable`, *optional*):
A function that calls every `callback_steps` steps during inference. The function is called with the
following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function is called. If not specified, the callback is called at
every step.
cross_attention_kwargs (`dict`, *optional*):
A kwargs dictionary that if specified is passed along to the [`AttentionProcessor`] as defined in
[`self.processor`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
Examples:
Returns:
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`:
If `return_dict` is `True`, [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] is returned,
otherwise a `tuple` is returned where the first element is a list with the generated images and the
second element is a list of `bool`s indicating whether the corresponding generated image contains
"not-safe-for-work" (nsfw) content.
"""
# 1. Check inputs. Raise error if not correct
self.check_inputs(
prompt,
rgb,
noise_level,
callback_steps,
negative_prompt,
prompt_embeds,
negative_prompt_embeds,
)
# 2. Define call parameters
if prompt is not None and isinstance(prompt, str):
batch_size = 1
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds.shape[0]
device = self._execution_device
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# 3. Encode input prompt
prompt_embeds, negative_prompt_embeds = self.encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
)
# For classifier free guidance, we need to do two forward passes.
# Here we concatenate the unconditional and text embeddings into a single batch
# to avoid doing two forward passes
if do_classifier_free_guidance:
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds])
# 4. Preprocess image
rgb, depth = self.image_processor.preprocess(rgb, depth, target_res=target_res)
rgb = rgb.to(dtype=prompt_embeds.dtype, device=device)
depth = depth.to(dtype=prompt_embeds.dtype, device=device)
# 5. set timesteps
self.scheduler.set_timesteps(num_inference_steps, device=device)
timesteps = self.scheduler.timesteps
# 6. Encode low resolutiom image to latent space
image = torch.cat([rgb, depth], axis=1)
latent_space_image = self.vae.encode(image).latent_dist.sample(generator)
latent_space_image *= self.vae.scaling_factor
noise_level = torch.tensor([noise_level], dtype=torch.long, device=device)
# noise_rgb = randn_tensor(rgb.shape, generator=generator, device=device, dtype=prompt_embeds.dtype)
# rgb = self.low_res_scheduler.add_noise(rgb, noise_rgb, noise_level)
# noise_depth = randn_tensor(depth.shape, generator=generator, device=device, dtype=prompt_embeds.dtype)
# depth = self.low_res_scheduler.add_noise(depth, noise_depth, noise_level)
batch_multiplier = 2 if do_classifier_free_guidance else 1
latent_space_image = torch.cat([latent_space_image] * batch_multiplier * num_images_per_prompt)
noise_level = torch.cat([noise_level] * latent_space_image.shape[0])
# 7. Prepare latent variables
height, width = latent_space_image.shape[2:]
num_channels_latents = self.vae.config.latent_channels
latents = self.prepare_latents(
batch_size * num_images_per_prompt,
num_channels_latents,
height,
width,
prompt_embeds.dtype,
device,
generator,
latents,
)
# 8. Check that sizes of image and latents match
num_channels_image = latent_space_image.shape[1]
if num_channels_latents + num_channels_image != self.unet.config.in_channels:
raise ValueError(
f"Incorrect configuration settings! The config of `pipeline.unet`: {self.unet.config} expects"
f" {self.unet.config.in_channels} but received `num_channels_latents`: {num_channels_latents} +"
f" `num_channels_image`: {num_channels_image} "
f" = {num_channels_latents+num_channels_image}. Please verify the config of"
" `pipeline.unet` or your `image` input."
)
# 9. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 10. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
# expand the latents if we are doing classifier free guidance
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
# concat latents, mask, masked_image_latents in the channel dimension
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
latent_model_input = torch.cat([latent_model_input, latent_space_image], dim=1)
# predict the noise residual
noise_pred = self.unet(
latent_model_input,
t,
encoder_hidden_states=prompt_embeds,
cross_attention_kwargs=cross_attention_kwargs,
class_labels=noise_level,
return_dict=False,
)[0]
# perform guidance
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
# compute the previous noisy sample x_t -> x_t-1
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0]
# call the callback, if provided
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0):
progress_bar.update()
if callback is not None and i % callback_steps == 0:
callback(i, t, latents)
if not output_type == "latent":
# make sure the VAE is in float32 mode, as it overflows in float16
needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast
if needs_upcasting:
self.upcast_vae()
latents = latents.to(next(iter(self.vae.post_quant_conv.parameters())).dtype)
image = self.vae.decode(latents / self.vae.scaling_factor, return_dict=False)[0]
# cast back to fp16 if needed
if needs_upcasting:
self.vae.to(dtype=torch.float16)
image, has_nsfw_concept = self.run_safety_checker(image, device, prompt_embeds.dtype)
else:
image = latents
has_nsfw_concept = None
if has_nsfw_concept is None:
do_denormalize = [True] * image.shape[0]
else:
do_denormalize = [not has_nsfw for has_nsfw in has_nsfw_concept]
rgb, depth = self.image_processor.postprocess(image, output_type=output_type, do_denormalize=do_denormalize)
# 11. Apply watermark
if output_type == "pil" and self.watermarker is not None:
rgb = self.watermarker.apply_watermark(rgb)
# Offload last model to CPU
if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None:
self.final_offload_hook.offload()
if not return_dict:
return ((rgb, depth), has_nsfw_concept)
return LDM3DPipelineOutput(rgb=rgb, depth=depth, nsfw_content_detected=has_nsfw_concept)
File diff suppressed because it is too large Load Diff
File diff suppressed because it is too large Load Diff
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import math
from typing import Dict, Optional
import torch
import torchvision.transforms.functional as FF
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from diffusers import StableDiffusionPipeline
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
from diffusers.schedulers import KarrasDiffusionSchedulers
from diffusers.utils import USE_PEFT_BACKEND
try:
from compel import Compel
except ImportError:
Compel = None
KCOMM = "ADDCOMM"
KBRK = "BREAK"
class RegionalPromptingStableDiffusionPipeline(StableDiffusionPipeline):
r"""
Args for Regional Prompting Pipeline:
rp_args:dict
Required
rp_args["mode"]: cols, rows, prompt, prompt-ex
for cols, rows mode
rp_args["div"]: ex) 1;1;1(Divide into 3 regions)
for prompt, prompt-ex mode
rp_args["th"]: ex) 0.5,0.5,0.6 (threshold for prompt mode)
Optional
rp_args["save_mask"]: True/False (save masks in prompt mode)
Pipeline for text-to-image generation using Stable Diffusion.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
text_encoder ([`CLIPTextModel`]):
Frozen text-encoder. Stable Diffusion uses the text portion of
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant.
tokenizer (`CLIPTokenizer`):
Tokenizer of class
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents.
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/CompVis/stable-diffusion-v1-4) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: KarrasDiffusionSchedulers,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
requires_safety_checker: bool = True,
):
super().__init__(
vae, text_encoder, tokenizer, unet, scheduler, safety_checker, feature_extractor, requires_safety_checker
)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
@torch.no_grad()
def __call__(
self,
prompt: str,
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: str = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
rp_args: Dict[str, str] = None,
):
active = KBRK in prompt[0] if type(prompt) == list else KBRK in prompt # noqa: E721
if negative_prompt is None:
negative_prompt = "" if type(prompt) == str else [""] * len(prompt) # noqa: E721
device = self._execution_device
regions = 0
self.power = int(rp_args["power"]) if "power" in rp_args else 1
prompts = prompt if type(prompt) == list else [prompt] # noqa: E721
n_prompts = negative_prompt if type(negative_prompt) == list else [negative_prompt] # noqa: E721
self.batch = batch = num_images_per_prompt * len(prompts)
all_prompts_cn, all_prompts_p = promptsmaker(prompts, num_images_per_prompt)
all_n_prompts_cn, _ = promptsmaker(n_prompts, num_images_per_prompt)
cn = len(all_prompts_cn) == len(all_n_prompts_cn)
if Compel:
compel = Compel(tokenizer=self.tokenizer, text_encoder=self.text_encoder)
def getcompelembs(prps):
embl = []
for prp in prps:
embl.append(compel.build_conditioning_tensor(prp))
return torch.cat(embl)
conds = getcompelembs(all_prompts_cn)
unconds = getcompelembs(all_n_prompts_cn) if cn else getcompelembs(n_prompts)
embs = getcompelembs(prompts)
n_embs = getcompelembs(n_prompts)
prompt = negative_prompt = None
else:
conds = self.encode_prompt(prompts, device, 1, True)[0]
unconds = (
self.encode_prompt(n_prompts, device, 1, True)[0]
if cn
else self.encode_prompt(all_n_prompts_cn, device, 1, True)[0]
)
embs = n_embs = None
if not active:
pcallback = None
mode = None
else:
if any(x in rp_args["mode"].upper() for x in ["COL", "ROW"]):
mode = "COL" if "COL" in rp_args["mode"].upper() else "ROW"
ocells, icells, regions = make_cells(rp_args["div"])
elif "PRO" in rp_args["mode"].upper():
regions = len(all_prompts_p[0])
mode = "PROMPT"
reset_attnmaps(self)
self.ex = "EX" in rp_args["mode"].upper()
self.target_tokens = target_tokens = tokendealer(self, all_prompts_p)
thresholds = [float(x) for x in rp_args["th"].split(",")]
orig_hw = (height, width)
revers = True
def pcallback(s_self, step: int, timestep: int, latents: torch.FloatTensor, selfs=None):
if "PRO" in mode: # in Prompt mode, make masks from sum of attension maps
self.step = step
if len(self.attnmaps_sizes) > 3:
self.history[step] = self.attnmaps.copy()
for hw in self.attnmaps_sizes:
allmasks = []
basemasks = [None] * batch
for tt, th in zip(target_tokens, thresholds):
for b in range(batch):
key = f"{tt}-{b}"
_, mask, _ = makepmask(self, self.attnmaps[key], hw[0], hw[1], th, step)
mask = mask.unsqueeze(0).unsqueeze(-1)
if self.ex:
allmasks[b::batch] = [x - mask for x in allmasks[b::batch]]
allmasks[b::batch] = [torch.where(x > 0, 1, 0) for x in allmasks[b::batch]]
allmasks.append(mask)
basemasks[b] = mask if basemasks[b] is None else basemasks[b] + mask
basemasks = [1 - mask for mask in basemasks]
basemasks = [torch.where(x > 0, 1, 0) for x in basemasks]
allmasks = basemasks + allmasks
self.attnmasks[hw] = torch.cat(allmasks)
self.maskready = True
return latents
def hook_forward(module):
# diffusers==0.23.2
def forward(
hidden_states: torch.FloatTensor,
encoder_hidden_states: Optional[torch.FloatTensor] = None,
attention_mask: Optional[torch.FloatTensor] = None,
temb: Optional[torch.FloatTensor] = None,
scale: float = 1.0,
) -> torch.Tensor:
attn = module
xshape = hidden_states.shape
self.hw = (h, w) = split_dims(xshape[1], *orig_hw)
if revers:
nx, px = hidden_states.chunk(2)
else:
px, nx = hidden_states.chunk(2)
if cn:
hidden_states = torch.cat([px for i in range(regions)] + [nx for i in range(regions)], 0)
encoder_hidden_states = torch.cat([conds] + [unconds])
else:
hidden_states = torch.cat([px for i in range(regions)] + [nx], 0)
encoder_hidden_states = torch.cat([conds] + [unconds])
residual = hidden_states
args = () if USE_PEFT_BACKEND else (scale,)
if attn.spatial_norm is not None:
hidden_states = attn.spatial_norm(hidden_states, temb)
input_ndim = hidden_states.ndim
if input_ndim == 4:
batch_size, channel, height, width = hidden_states.shape
hidden_states = hidden_states.view(batch_size, channel, height * width).transpose(1, 2)
batch_size, sequence_length, _ = (
hidden_states.shape if encoder_hidden_states is None else encoder_hidden_states.shape
)
if attention_mask is not None:
attention_mask = attn.prepare_attention_mask(attention_mask, sequence_length, batch_size)
attention_mask = attention_mask.view(batch_size, attn.heads, -1, attention_mask.shape[-1])
if attn.group_norm is not None:
hidden_states = attn.group_norm(hidden_states.transpose(1, 2)).transpose(1, 2)
args = () if USE_PEFT_BACKEND else (scale,)
query = attn.to_q(hidden_states, *args)
if encoder_hidden_states is None:
encoder_hidden_states = hidden_states
elif attn.norm_cross:
encoder_hidden_states = attn.norm_encoder_hidden_states(encoder_hidden_states)
key = attn.to_k(encoder_hidden_states, *args)
value = attn.to_v(encoder_hidden_states, *args)
inner_dim = key.shape[-1]
head_dim = inner_dim // attn.heads
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
key = key.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
value = value.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
# the output of sdp = (batch, num_heads, seq_len, head_dim)
# TODO: add support for attn.scale when we move to Torch 2.1
hidden_states = scaled_dot_product_attention(
self,
query,
key,
value,
attn_mask=attention_mask,
dropout_p=0.0,
is_causal=False,
getattn="PRO" in mode,
)
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
hidden_states = hidden_states.to(query.dtype)
# linear proj
hidden_states = attn.to_out[0](hidden_states, *args)
# dropout
hidden_states = attn.to_out[1](hidden_states)
if input_ndim == 4:
hidden_states = hidden_states.transpose(-1, -2).reshape(batch_size, channel, height, width)
if attn.residual_connection:
hidden_states = hidden_states + residual
hidden_states = hidden_states / attn.rescale_output_factor
#### Regional Prompting Col/Row mode
if any(x in mode for x in ["COL", "ROW"]):
reshaped = hidden_states.reshape(hidden_states.size()[0], h, w, hidden_states.size()[2])
center = reshaped.shape[0] // 2
px = reshaped[0:center] if cn else reshaped[0:-batch]
nx = reshaped[center:] if cn else reshaped[-batch:]
outs = [px, nx] if cn else [px]
for out in outs:
c = 0
for i, ocell in enumerate(ocells):
for icell in icells[i]:
if "ROW" in mode:
out[
0:batch,
int(h * ocell[0]) : int(h * ocell[1]),
int(w * icell[0]) : int(w * icell[1]),
:,
] = out[
c * batch : (c + 1) * batch,
int(h * ocell[0]) : int(h * ocell[1]),
int(w * icell[0]) : int(w * icell[1]),
:,
]
else:
out[
0:batch,
int(h * icell[0]) : int(h * icell[1]),
int(w * ocell[0]) : int(w * ocell[1]),
:,
] = out[
c * batch : (c + 1) * batch,
int(h * icell[0]) : int(h * icell[1]),
int(w * ocell[0]) : int(w * ocell[1]),
:,
]
c += 1
px, nx = (px[0:batch], nx[0:batch]) if cn else (px[0:batch], nx)
hidden_states = torch.cat([nx, px], 0) if revers else torch.cat([px, nx], 0)
hidden_states = hidden_states.reshape(xshape)
#### Regional Prompting Prompt mode
elif "PRO" in mode:
center = reshaped.shape[0] // 2
px = reshaped[0:center] if cn else reshaped[0:-batch]
nx = reshaped[center:] if cn else reshaped[-batch:]
if (h, w) in self.attnmasks and self.maskready:
def mask(input):
out = torch.multiply(input, self.attnmasks[(h, w)])
for b in range(batch):
for r in range(1, regions):
out[b] = out[b] + out[r * batch + b]
return out
px, nx = (mask(px), mask(nx)) if cn else (mask(px), nx)
px, nx = (px[0:batch], nx[0:batch]) if cn else (px[0:batch], nx)
hidden_states = torch.cat([nx, px], 0) if revers else torch.cat([px, nx], 0)
return hidden_states
return forward
def hook_forwards(root_module: torch.nn.Module):
for name, module in root_module.named_modules():
if "attn2" in name and module.__class__.__name__ == "Attention":
module.forward = hook_forward(module)
hook_forwards(self.unet)
output = StableDiffusionPipeline(**self.components)(
prompt=prompt,
prompt_embeds=embs,
negative_prompt=negative_prompt,
negative_prompt_embeds=n_embs,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback_on_step_end=pcallback,
)
if "save_mask" in rp_args:
save_mask = rp_args["save_mask"]
else:
save_mask = False
if mode == "PROMPT" and save_mask:
saveattnmaps(self, output, height, width, thresholds, num_inference_steps // 2, regions)
return output
### Make prompt list for each regions
def promptsmaker(prompts, batch):
out_p = []
plen = len(prompts)
for prompt in prompts:
add = ""
if KCOMM in prompt:
add, prompt = prompt.split(KCOMM)
add = add + " "
prompts = prompt.split(KBRK)
out_p.append([add + p for p in prompts])
out = [None] * batch * len(out_p[0]) * len(out_p)
for p, prs in enumerate(out_p): # inputs prompts
for r, pr in enumerate(prs): # prompts for regions
start = (p + r * plen) * batch
out[start : start + batch] = [pr] * batch # P1R1B1,P1R1B2...,P1R2B1,P1R2B2...,P2R1B1...
return out, out_p
### make regions from ratios
### ";" makes outercells, "," makes inner cells
def make_cells(ratios):
if ";" not in ratios and "," in ratios:
ratios = ratios.replace(",", ";")
ratios = ratios.split(";")
ratios = [inratios.split(",") for inratios in ratios]
icells = []
ocells = []
def startend(cells, array):
current_start = 0
array = [float(x) for x in array]
for value in array:
end = current_start + (value / sum(array))
cells.append([current_start, end])
current_start = end
startend(ocells, [r[0] for r in ratios])
for inratios in ratios:
if 2 > len(inratios):
icells.append([[0, 1]])
else:
add = []
startend(add, inratios[1:])
icells.append(add)
return ocells, icells, sum(len(cell) for cell in icells)
def make_emblist(self, prompts):
with torch.no_grad():
tokens = self.tokenizer(
prompts, max_length=self.tokenizer.model_max_length, padding=True, truncation=True, return_tensors="pt"
).input_ids.to(self.device)
embs = self.text_encoder(tokens, output_hidden_states=True).last_hidden_state.to(self.device, dtype=self.dtype)
return embs
def split_dims(xs, height, width):
xs = xs
def repeat_div(x, y):
while y > 0:
x = math.ceil(x / 2)
y = y - 1
return x
scale = math.ceil(math.log2(math.sqrt(height * width / xs)))
dsh = repeat_div(height, scale)
dsw = repeat_div(width, scale)
return dsh, dsw
##### for prompt mode
def get_attn_maps(self, attn):
height, width = self.hw
target_tokens = self.target_tokens
if (height, width) not in self.attnmaps_sizes:
self.attnmaps_sizes.append((height, width))
for b in range(self.batch):
for t in target_tokens:
power = self.power
add = attn[b, :, :, t[0] : t[0] + len(t)] ** (power) * (self.attnmaps_sizes.index((height, width)) + 1)
add = torch.sum(add, dim=2)
key = f"{t}-{b}"
if key not in self.attnmaps:
self.attnmaps[key] = add
else:
if self.attnmaps[key].shape[1] != add.shape[1]:
add = add.view(8, height, width)
add = FF.resize(add, self.attnmaps_sizes[0], antialias=None)
add = add.reshape_as(self.attnmaps[key])
self.attnmaps[key] = self.attnmaps[key] + add
def reset_attnmaps(self): # init parameters in every batch
self.step = 0
self.attnmaps = {} # maked from attention maps
self.attnmaps_sizes = [] # height,width set of u-net blocks
self.attnmasks = {} # maked from attnmaps for regions
self.maskready = False
self.history = {}
def saveattnmaps(self, output, h, w, th, step, regions):
masks = []
for i, mask in enumerate(self.history[step].values()):
img, _, mask = makepmask(self, mask, h, w, th[i % len(th)], step)
if self.ex:
masks = [x - mask for x in masks]
masks.append(mask)
if len(masks) == regions - 1:
output.images.extend([FF.to_pil_image(mask) for mask in masks])
masks = []
else:
output.images.append(img)
def makepmask(
self, mask, h, w, th, step
): # make masks from attention cache return [for preview, for attention, for Latent]
th = th - step * 0.005
if 0.05 >= th:
th = 0.05
mask = torch.mean(mask, dim=0)
mask = mask / mask.max().item()
mask = torch.where(mask > th, 1, 0)
mask = mask.float()
mask = mask.view(1, *self.attnmaps_sizes[0])
img = FF.to_pil_image(mask)
img = img.resize((w, h))
mask = FF.resize(mask, (h, w), interpolation=FF.InterpolationMode.NEAREST, antialias=None)
lmask = mask
mask = mask.reshape(h * w)
mask = torch.where(mask > 0.1, 1, 0)
return img, mask, lmask
def tokendealer(self, all_prompts):
for prompts in all_prompts:
targets = [p.split(",")[-1] for p in prompts[1:]]
tt = []
for target in targets:
ptokens = (
self.tokenizer(
prompts,
max_length=self.tokenizer.model_max_length,
padding=True,
truncation=True,
return_tensors="pt",
).input_ids
)[0]
ttokens = (
self.tokenizer(
target,
max_length=self.tokenizer.model_max_length,
padding=True,
truncation=True,
return_tensors="pt",
).input_ids
)[0]
tlist = []
for t in range(ttokens.shape[0] - 2):
for p in range(ptokens.shape[0]):
if ttokens[t + 1] == ptokens[p]:
tlist.append(p)
if tlist != []:
tt.append(tlist)
return tt
def scaled_dot_product_attention(
self, query, key, value, attn_mask=None, dropout_p=0.0, is_causal=False, scale=None, getattn=False
) -> torch.Tensor:
# Efficient implementation equivalent to the following:
L, S = query.size(-2), key.size(-2)
scale_factor = 1 / math.sqrt(query.size(-1)) if scale is None else scale
attn_bias = torch.zeros(L, S, dtype=query.dtype, device=self.device)
if is_causal:
assert attn_mask is None
temp_mask = torch.ones(L, S, dtype=torch.bool).tril(diagonal=0)
attn_bias.masked_fill_(temp_mask.logical_not(), float("-inf"))
attn_bias.to(query.dtype)
if attn_mask is not None:
if attn_mask.dtype == torch.bool:
attn_mask.masked_fill_(attn_mask.logical_not(), float("-inf"))
else:
attn_bias += attn_mask
attn_weight = query @ key.transpose(-2, -1) * scale_factor
attn_weight += attn_bias
attn_weight = torch.softmax(attn_weight, dim=-1)
if getattn:
get_attn_maps(self, attn_weight)
attn_weight = torch.dropout(attn_weight, dropout_p, train=True)
return attn_weight @ value
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import math
import tempfile
from typing import List, Optional
import numpy as np
import PIL.Image
import torch
from accelerate import Accelerator
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
from diffusers import AutoencoderKL, DiffusionPipeline, DPMSolverMultistepScheduler, UNet2DConditionModel
from diffusers.loaders import AttnProcsLayers, LoraLoaderMixin
from diffusers.models.attention_processor import (
AttnAddedKVProcessor,
AttnAddedKVProcessor2_0,
LoRAAttnAddedKVProcessor,
LoRAAttnProcessor,
LoRAAttnProcessor2_0,
SlicedAttnAddedKVProcessor,
)
from diffusers.optimization import get_scheduler
class SdeDragPipeline(DiffusionPipeline):
r"""
Pipeline for image drag-and-drop editing using stochastic differential equations: https://arxiv.org/abs/2311.01410.
Please refer to the [official repository](https://github.com/ML-GSAI/SDE-Drag) for more information.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
text_encoder ([`CLIPTextModel`]):
Frozen text-encoder. Stable Diffusion uses the text portion of
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant.
tokenizer (`CLIPTokenizer`):
Tokenizer of class
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents.
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Please use
[`DDIMScheduler`].
"""
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: DPMSolverMultistepScheduler,
):
super().__init__()
self.register_modules(vae=vae, text_encoder=text_encoder, tokenizer=tokenizer, unet=unet, scheduler=scheduler)
@torch.no_grad()
def __call__(
self,
prompt: str,
image: PIL.Image.Image,
mask_image: PIL.Image.Image,
source_points: List[List[int]],
target_points: List[List[int]],
t0: Optional[float] = 0.6,
steps: Optional[int] = 200,
step_size: Optional[int] = 2,
image_scale: Optional[float] = 0.3,
adapt_radius: Optional[int] = 5,
min_lora_scale: Optional[float] = 0.5,
generator: Optional[torch.Generator] = None,
):
r"""
Function invoked when calling the pipeline for image editing.
Args:
prompt (`str`, *required*):
The prompt to guide the image editing.
image (`PIL.Image.Image`, *required*):
Which will be edited, parts of the image will be masked out with `mask_image` and edited
according to `prompt`.
mask_image (`PIL.Image.Image`, *required*):
To mask `image`. White pixels in the mask will be edited, while black pixels will be preserved.
source_points (`List[List[int]]`, *required*):
Used to mark the starting positions of drag editing in the image, with each pixel represented as a
`List[int]` of length 2.
target_points (`List[List[int]]`, *required*):
Used to mark the target positions of drag editing in the image, with each pixel represented as a
`List[int]` of length 2.
t0 (`float`, *optional*, defaults to 0.6):
The time parameter. Higher t0 improves the fidelity while lowering the faithfulness of the edited images
and vice versa.
steps (`int`, *optional*, defaults to 200):
The number of sampling iterations.
step_size (`int`, *optional*, defaults to 2):
The drag diatance of each drag step.
image_scale (`float`, *optional*, defaults to 0.3):
To avoid duplicating the content, use image_scale to perturbs the source.
adapt_radius (`int`, *optional*, defaults to 5):
The size of the region for copy and paste operations during each step of the drag process.
min_lora_scale (`float`, *optional*, defaults to 0.5):
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
min_lora_scale specifies the minimum LoRA scale during the image drag-editing process.
generator ('torch.Generator', *optional*, defaults to None):
To make generation deterministic(https://pytorch.org/docs/stable/generated/torch.Generator.html).
Examples:
```py
>>> import PIL
>>> import torch
>>> from diffusers import DDIMScheduler, DiffusionPipeline
>>> # Load the pipeline
>>> model_path = "runwayml/stable-diffusion-v1-5"
>>> scheduler = DDIMScheduler.from_pretrained(model_path, subfolder="scheduler")
>>> pipe = DiffusionPipeline.from_pretrained(model_path, scheduler=scheduler, custom_pipeline="sde_drag")
>>> pipe.to('cuda')
>>> # To save GPU memory, torch.float16 can be used, but it may compromise image quality.
>>> # If not training LoRA, please avoid using torch.float16
>>> # pipe.to(torch.float16)
>>> # Provide prompt, image, mask image, and the starting and target points for drag editing.
>>> prompt = "prompt of the image"
>>> image = PIL.Image.open('/path/to/image')
>>> mask_image = PIL.Image.open('/path/to/mask_image')
>>> source_points = [[123, 456]]
>>> target_points = [[234, 567]]
>>> # train_lora is optional, and in most cases, using train_lora can better preserve consistency with the original image.
>>> pipe.train_lora(prompt, image)
>>> output = pipe(prompt, image, mask_image, source_points, target_points)
>>> output_image = PIL.Image.fromarray(output)
>>> output_image.save("./output.png")
```
"""
self.scheduler.set_timesteps(steps)
noise_scale = (1 - image_scale**2) ** (0.5)
text_embeddings = self._get_text_embed(prompt)
uncond_embeddings = self._get_text_embed([""])
text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
latent = self._get_img_latent(image)
mask = mask_image.resize((latent.shape[3], latent.shape[2]))
mask = torch.tensor(np.array(mask))
mask = mask.unsqueeze(0).expand_as(latent).to(self.device)
source_points = torch.tensor(source_points).div(torch.tensor([8]), rounding_mode="trunc")
target_points = torch.tensor(target_points).div(torch.tensor([8]), rounding_mode="trunc")
distance = target_points - source_points
distance_norm_max = torch.norm(distance.float(), dim=1, keepdim=True).max()
if distance_norm_max <= step_size:
drag_num = 1
else:
drag_num = distance_norm_max.div(torch.tensor([step_size]), rounding_mode="trunc")
if (distance_norm_max / drag_num - step_size).abs() > (
distance_norm_max / (drag_num + 1) - step_size
).abs():
drag_num += 1
latents = []
for i in tqdm(range(int(drag_num)), desc="SDE Drag"):
source_new = source_points + (i / drag_num * distance).to(torch.int)
target_new = source_points + ((i + 1) / drag_num * distance).to(torch.int)
latent, noises, hook_latents, lora_scales, cfg_scales = self._forward(
latent, steps, t0, min_lora_scale, text_embeddings, generator
)
latent = self._copy_and_paste(
latent,
source_new,
target_new,
adapt_radius,
latent.shape[2] - 1,
latent.shape[3] - 1,
image_scale,
noise_scale,
generator,
)
latent = self._backward(
latent, mask, steps, t0, noises, hook_latents, lora_scales, cfg_scales, text_embeddings, generator
)
latents.append(latent)
result_image = 1 / 0.18215 * latents[-1]
with torch.no_grad():
result_image = self.vae.decode(result_image).sample
result_image = (result_image / 2 + 0.5).clamp(0, 1)
result_image = result_image.cpu().permute(0, 2, 3, 1).numpy()[0]
result_image = (result_image * 255).astype(np.uint8)
return result_image
def train_lora(self, prompt, image, lora_step=100, lora_rank=16, generator=None):
accelerator = Accelerator(gradient_accumulation_steps=1, mixed_precision="fp16")
self.vae.requires_grad_(False)
self.text_encoder.requires_grad_(False)
self.unet.requires_grad_(False)
unet_lora_attn_procs = {}
for name, attn_processor in self.unet.attn_processors.items():
cross_attention_dim = None if name.endswith("attn1.processor") else self.unet.config.cross_attention_dim
if name.startswith("mid_block"):
hidden_size = self.unet.config.block_out_channels[-1]
elif name.startswith("up_blocks"):
block_id = int(name[len("up_blocks.")])
hidden_size = list(reversed(self.unet.config.block_out_channels))[block_id]
elif name.startswith("down_blocks"):
block_id = int(name[len("down_blocks.")])
hidden_size = self.unet.config.block_out_channels[block_id]
else:
raise NotImplementedError("name must start with up_blocks, mid_blocks, or down_blocks")
if isinstance(attn_processor, (AttnAddedKVProcessor, SlicedAttnAddedKVProcessor, AttnAddedKVProcessor2_0)):
lora_attn_processor_class = LoRAAttnAddedKVProcessor
else:
lora_attn_processor_class = (
LoRAAttnProcessor2_0
if hasattr(torch.nn.functional, "scaled_dot_product_attention")
else LoRAAttnProcessor
)
unet_lora_attn_procs[name] = lora_attn_processor_class(
hidden_size=hidden_size, cross_attention_dim=cross_attention_dim, rank=lora_rank
)
self.unet.set_attn_processor(unet_lora_attn_procs)
unet_lora_layers = AttnProcsLayers(self.unet.attn_processors)
params_to_optimize = unet_lora_layers.parameters()
optimizer = torch.optim.AdamW(
params_to_optimize,
lr=2e-4,
betas=(0.9, 0.999),
weight_decay=1e-2,
eps=1e-08,
)
lr_scheduler = get_scheduler(
"constant",
optimizer=optimizer,
num_warmup_steps=0,
num_training_steps=lora_step,
num_cycles=1,
power=1.0,
)
unet_lora_layers = accelerator.prepare_model(unet_lora_layers)
optimizer = accelerator.prepare_optimizer(optimizer)
lr_scheduler = accelerator.prepare_scheduler(lr_scheduler)
with torch.no_grad():
text_inputs = self._tokenize_prompt(prompt, tokenizer_max_length=None)
text_embedding = self._encode_prompt(
text_inputs.input_ids, text_inputs.attention_mask, text_encoder_use_attention_mask=False
)
image_transforms = transforms.Compose(
[
transforms.ToTensor(),
transforms.Normalize([0.5], [0.5]),
]
)
image = image_transforms(image).to(self.device, dtype=self.vae.dtype)
image = image.unsqueeze(dim=0)
latents_dist = self.vae.encode(image).latent_dist
for _ in tqdm(range(lora_step), desc="Train LoRA"):
self.unet.train()
model_input = latents_dist.sample() * self.vae.config.scaling_factor
# Sample noise that we'll add to the latents
noise = torch.randn(
model_input.size(),
dtype=model_input.dtype,
layout=model_input.layout,
device=model_input.device,
generator=generator,
)
bsz, channels, height, width = model_input.shape
# Sample a random timestep for each image
timesteps = torch.randint(
0, self.scheduler.config.num_train_timesteps, (bsz,), device=model_input.device, generator=generator
)
timesteps = timesteps.long()
# Add noise to the model input according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_model_input = self.scheduler.add_noise(model_input, noise, timesteps)
# Predict the noise residual
model_pred = self.unet(noisy_model_input, timesteps, text_embedding).sample
# Get the target for loss depending on the prediction type
if self.scheduler.config.prediction_type == "epsilon":
target = noise
elif self.scheduler.config.prediction_type == "v_prediction":
target = self.scheduler.get_velocity(model_input, noise, timesteps)
else:
raise ValueError(f"Unknown prediction type {self.scheduler.config.prediction_type}")
loss = torch.nn.functional.mse_loss(model_pred.float(), target.float(), reduction="mean")
accelerator.backward(loss)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
with tempfile.TemporaryDirectory() as save_lora_dir:
LoraLoaderMixin.save_lora_weights(
save_directory=save_lora_dir,
unet_lora_layers=unet_lora_layers,
text_encoder_lora_layers=None,
)
self.unet.load_attn_procs(save_lora_dir)
def _tokenize_prompt(self, prompt, tokenizer_max_length=None):
if tokenizer_max_length is not None:
max_length = tokenizer_max_length
else:
max_length = self.tokenizer.model_max_length
text_inputs = self.tokenizer(
prompt,
truncation=True,
padding="max_length",
max_length=max_length,
return_tensors="pt",
)
return text_inputs
def _encode_prompt(self, input_ids, attention_mask, text_encoder_use_attention_mask=False):
text_input_ids = input_ids.to(self.device)
if text_encoder_use_attention_mask:
attention_mask = attention_mask.to(self.device)
else:
attention_mask = None
prompt_embeds = self.text_encoder(
text_input_ids,
attention_mask=attention_mask,
)
prompt_embeds = prompt_embeds[0]
return prompt_embeds
@torch.no_grad()
def _get_text_embed(self, prompt):
text_input = self.tokenizer(
prompt,
padding="max_length",
max_length=self.tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_embeddings = self.text_encoder(text_input.input_ids.to(self.device))[0]
return text_embeddings
def _copy_and_paste(
self, latent, source_new, target_new, adapt_radius, max_height, max_width, image_scale, noise_scale, generator
):
def adaption_r(source, target, adapt_radius, max_height, max_width):
r_x_lower = min(adapt_radius, source[0], target[0])
r_x_upper = min(adapt_radius, max_width - source[0], max_width - target[0])
r_y_lower = min(adapt_radius, source[1], target[1])
r_y_upper = min(adapt_radius, max_height - source[1], max_height - target[1])
return r_x_lower, r_x_upper, r_y_lower, r_y_upper
for source_, target_ in zip(source_new, target_new):
r_x_lower, r_x_upper, r_y_lower, r_y_upper = adaption_r(
source_, target_, adapt_radius, max_height, max_width
)
source_feature = latent[
:, :, source_[1] - r_y_lower : source_[1] + r_y_upper, source_[0] - r_x_lower : source_[0] + r_x_upper
].clone()
latent[
:, :, source_[1] - r_y_lower : source_[1] + r_y_upper, source_[0] - r_x_lower : source_[0] + r_x_upper
] = image_scale * source_feature + noise_scale * torch.randn(
latent.shape[0],
4,
r_y_lower + r_y_upper,
r_x_lower + r_x_upper,
device=self.device,
generator=generator,
)
latent[
:, :, target_[1] - r_y_lower : target_[1] + r_y_upper, target_[0] - r_x_lower : target_[0] + r_x_upper
] = source_feature * 1.1
return latent
@torch.no_grad()
def _get_img_latent(self, image, height=None, weight=None):
data = image.convert("RGB")
if height is not None:
data = data.resize((weight, height))
transform = transforms.ToTensor()
data = transform(data).unsqueeze(0)
data = (data * 2.0) - 1.0
data = data.to(self.device, dtype=self.vae.dtype)
latent = self.vae.encode(data).latent_dist.sample()
latent = 0.18215 * latent
return latent
@torch.no_grad()
def _get_eps(self, latent, timestep, guidance_scale, text_embeddings, lora_scale=None):
latent_model_input = torch.cat([latent] * 2) if guidance_scale > 1.0 else latent
text_embeddings = text_embeddings if guidance_scale > 1.0 else text_embeddings.chunk(2)[1]
cross_attention_kwargs = None if lora_scale is None else {"scale": lora_scale}
with torch.no_grad():
noise_pred = self.unet(
latent_model_input,
timestep,
encoder_hidden_states=text_embeddings,
cross_attention_kwargs=cross_attention_kwargs,
).sample
if guidance_scale > 1.0:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
elif guidance_scale == 1.0:
noise_pred_text = noise_pred
noise_pred_uncond = 0.0
else:
raise NotImplementedError(guidance_scale)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
return noise_pred
def _forward_sde(
self, timestep, sample, guidance_scale, text_embeddings, steps, eta=1.0, lora_scale=None, generator=None
):
num_train_timesteps = len(self.scheduler)
alphas_cumprod = self.scheduler.alphas_cumprod
initial_alpha_cumprod = torch.tensor(1.0)
prev_timestep = timestep + num_train_timesteps // steps
alpha_prod_t = alphas_cumprod[timestep] if timestep >= 0 else initial_alpha_cumprod
alpha_prod_t_prev = alphas_cumprod[prev_timestep]
beta_prod_t_prev = 1 - alpha_prod_t_prev
x_prev = (alpha_prod_t_prev / alpha_prod_t) ** (0.5) * sample + (1 - alpha_prod_t_prev / alpha_prod_t) ** (
0.5
) * torch.randn(
sample.size(), dtype=sample.dtype, layout=sample.layout, device=self.device, generator=generator
)
eps = self._get_eps(x_prev, prev_timestep, guidance_scale, text_embeddings, lora_scale)
sigma_t_prev = (
eta
* (1 - alpha_prod_t) ** (0.5)
* (1 - alpha_prod_t_prev / (1 - alpha_prod_t_prev) * (1 - alpha_prod_t) / alpha_prod_t) ** (0.5)
)
pred_original_sample = (x_prev - beta_prod_t_prev ** (0.5) * eps) / alpha_prod_t_prev ** (0.5)
pred_sample_direction_coeff = (1 - alpha_prod_t - sigma_t_prev**2) ** (0.5)
noise = (
sample - alpha_prod_t ** (0.5) * pred_original_sample - pred_sample_direction_coeff * eps
) / sigma_t_prev
return x_prev, noise
def _sample(
self,
timestep,
sample,
guidance_scale,
text_embeddings,
steps,
sde=False,
noise=None,
eta=1.0,
lora_scale=None,
generator=None,
):
num_train_timesteps = len(self.scheduler)
alphas_cumprod = self.scheduler.alphas_cumprod
final_alpha_cumprod = torch.tensor(1.0)
eps = self._get_eps(sample, timestep, guidance_scale, text_embeddings, lora_scale)
prev_timestep = timestep - num_train_timesteps // steps
alpha_prod_t = alphas_cumprod[timestep]
alpha_prod_t_prev = alphas_cumprod[prev_timestep] if prev_timestep >= 0 else final_alpha_cumprod
beta_prod_t = 1 - alpha_prod_t
sigma_t = (
eta
* ((1 - alpha_prod_t_prev) / (1 - alpha_prod_t)) ** (0.5)
* (1 - alpha_prod_t / alpha_prod_t_prev) ** (0.5)
if sde
else 0
)
pred_original_sample = (sample - beta_prod_t ** (0.5) * eps) / alpha_prod_t ** (0.5)
pred_sample_direction_coeff = (1 - alpha_prod_t_prev - sigma_t**2) ** (0.5)
noise = (
torch.randn(
sample.size(), dtype=sample.dtype, layout=sample.layout, device=self.device, generator=generator
)
if noise is None
else noise
)
latent = (
alpha_prod_t_prev ** (0.5) * pred_original_sample + pred_sample_direction_coeff * eps + sigma_t * noise
)
return latent
def _forward(self, latent, steps, t0, lora_scale_min, text_embeddings, generator):
def scale_schedule(begin, end, n, length, type="linear"):
if type == "constant":
return end
elif type == "linear":
return begin + (end - begin) * n / length
elif type == "cos":
factor = (1 - math.cos(n * math.pi / length)) / 2
return (1 - factor) * begin + factor * end
else:
raise NotImplementedError(type)
noises = []
latents = []
lora_scales = []
cfg_scales = []
latents.append(latent)
t0 = int(t0 * steps)
t_begin = steps - t0
length = len(self.scheduler.timesteps[t_begin - 1 : -1]) - 1
index = 1
for t in self.scheduler.timesteps[t_begin:].flip(dims=[0]):
lora_scale = scale_schedule(1, lora_scale_min, index, length, type="cos")
cfg_scale = scale_schedule(1, 3.0, index, length, type="linear")
latent, noise = self._forward_sde(
t, latent, cfg_scale, text_embeddings, steps, lora_scale=lora_scale, generator=generator
)
noises.append(noise)
latents.append(latent)
lora_scales.append(lora_scale)
cfg_scales.append(cfg_scale)
index += 1
return latent, noises, latents, lora_scales, cfg_scales
def _backward(
self, latent, mask, steps, t0, noises, hook_latents, lora_scales, cfg_scales, text_embeddings, generator
):
t0 = int(t0 * steps)
t_begin = steps - t0
hook_latent = hook_latents.pop()
latent = torch.where(mask > 128, latent, hook_latent)
for t in self.scheduler.timesteps[t_begin - 1 : -1]:
latent = self._sample(
t,
latent,
cfg_scales.pop(),
text_embeddings,
steps,
sde=True,
noise=noises.pop(),
lora_scale=lora_scales.pop(),
generator=generator,
)
hook_latent = hook_latents.pop()
latent = torch.where(mask > 128, latent, hook_latent)
return latent
+2 -2
View File
@@ -21,7 +21,7 @@ from packaging import version
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from diffusers.configuration_utils import FrozenDict
from diffusers.loaders import TextualInversionLoaderMixin
from diffusers.loaders import LoraLoaderMixin, TextualInversionLoaderMixin
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.pipeline_utils import DiffusionPipeline
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
@@ -62,7 +62,7 @@ EXAMPLE_DOC_STRING = """
"""
class StableDiffusionIPEXPipeline(DiffusionPipeline, TextualInversionLoaderMixin):
class StableDiffusionIPEXPipeline(DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin):
r"""
Pipeline for text-to-image generation using Stable Diffusion on IPEX.
@@ -28,6 +28,7 @@ import PIL.Image
import tensorrt as trt
import torch
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
@@ -41,7 +42,7 @@ from polygraphy.backend.trt import (
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
@@ -50,7 +51,7 @@ from diffusers.pipelines.stable_diffusion import (
StableDiffusionSafetyChecker,
)
from diffusers.schedulers import DDIMScheduler
from diffusers.utils import DIFFUSERS_CACHE, logging
from diffusers.utils import logging
"""
@@ -709,6 +710,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
scheduler: DDIMScheduler,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
image_encoder: CLIPVisionModelWithProjection = None,
requires_safety_checker: bool = True,
stages=["clip", "unet", "vae", "vae_encoder"],
image_height: int = 512,
@@ -724,7 +726,15 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
timing_cache: str = "timing_cache",
):
super().__init__(
vae, text_encoder, tokenizer, unet, scheduler, safety_checker, feature_extractor, requires_safety_checker
vae,
text_encoder,
tokenizer,
unet,
scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
@@ -769,12 +779,13 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
self.models["vae_encoder"] = make_VAEEncoder(self.vae, **models_args)
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
cache_dir = kwargs.pop("cache_dir", DIFFUSERS_CACHE)
cache_dir = kwargs.pop("cache_dir", None)
resume_download = kwargs.pop("resume_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", False)
use_auth_token = kwargs.pop("use_auth_token", None)
token = kwargs.pop("token", None)
revision = kwargs.pop("revision", None)
cls.cached_folder = (
@@ -786,7 +797,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(StableDiffusionImg2ImgPipeline):
resume_download=resume_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
token=token,
revision=revision,
)
)
@@ -28,6 +28,7 @@ import PIL.Image
import tensorrt as trt
import torch
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
@@ -41,7 +42,7 @@ from polygraphy.backend.trt import (
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
@@ -51,7 +52,7 @@ from diffusers.pipelines.stable_diffusion import (
)
from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_inpaint import prepare_mask_and_masked_image
from diffusers.schedulers import DDIMScheduler
from diffusers.utils import DIFFUSERS_CACHE, logging
from diffusers.utils import logging
"""
@@ -710,6 +711,7 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
scheduler: DDIMScheduler,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
image_encoder: CLIPVisionModelWithProjection = None,
requires_safety_checker: bool = True,
stages=["clip", "unet", "vae", "vae_encoder"],
image_height: int = 512,
@@ -725,7 +727,15 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
timing_cache: str = "timing_cache",
):
super().__init__(
vae, text_encoder, tokenizer, unet, scheduler, safety_checker, feature_extractor, requires_safety_checker
vae,
text_encoder,
tokenizer,
unet,
scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
@@ -770,12 +780,13 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
self.models["vae_encoder"] = make_VAEEncoder(self.vae, **models_args)
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
cache_dir = kwargs.pop("cache_dir", DIFFUSERS_CACHE)
cache_dir = kwargs.pop("cache_dir", None)
resume_download = kwargs.pop("resume_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", False)
use_auth_token = kwargs.pop("use_auth_token", None)
token = kwargs.pop("token", None)
revision = kwargs.pop("revision", None)
cls.cached_folder = (
@@ -787,7 +798,7 @@ class TensorRTStableDiffusionInpaintPipeline(StableDiffusionInpaintPipeline):
resume_download=resume_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
token=token,
revision=revision,
)
)
@@ -27,6 +27,7 @@ import onnx_graphsurgeon as gs
import tensorrt as trt
import torch
from huggingface_hub import snapshot_download
from huggingface_hub.utils import validate_hf_hub_args
from onnx import shape_inference
from polygraphy import cuda
from polygraphy.backend.common import bytes_from_path
@@ -40,7 +41,7 @@ from polygraphy.backend.trt import (
save_engine,
)
from polygraphy.backend.trt import util as trt_util
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer, CLIPVisionModelWithProjection
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipelines.stable_diffusion import (
@@ -49,7 +50,7 @@ from diffusers.pipelines.stable_diffusion import (
StableDiffusionSafetyChecker,
)
from diffusers.schedulers import DDIMScheduler
from diffusers.utils import DIFFUSERS_CACHE, logging
from diffusers.utils import logging
"""
@@ -624,6 +625,7 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
scheduler: DDIMScheduler,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
image_encoder: CLIPVisionModelWithProjection = None,
requires_safety_checker: bool = True,
stages=["clip", "unet", "vae"],
image_height: int = 768,
@@ -639,7 +641,15 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
timing_cache: str = "timing_cache",
):
super().__init__(
vae, text_encoder, tokenizer, unet, scheduler, safety_checker, feature_extractor, requires_safety_checker
vae,
text_encoder,
tokenizer,
unet,
scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
requires_safety_checker=requires_safety_checker,
)
self.vae.forward = self.vae.decode
@@ -682,12 +692,13 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
self.models["vae"] = make_VAE(self.vae, **models_args)
@classmethod
@validate_hf_hub_args
def set_cached_folder(cls, pretrained_model_name_or_path: Optional[Union[str, os.PathLike]], **kwargs):
cache_dir = kwargs.pop("cache_dir", DIFFUSERS_CACHE)
cache_dir = kwargs.pop("cache_dir", None)
resume_download = kwargs.pop("resume_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", False)
use_auth_token = kwargs.pop("use_auth_token", None)
token = kwargs.pop("token", None)
revision = kwargs.pop("revision", None)
cls.cached_folder = (
@@ -699,7 +710,7 @@ class TensorRTStableDiffusionPipeline(StableDiffusionPipeline):
resume_download=resume_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
token=token,
revision=revision,
)
)
+22 -14
View File
@@ -1,6 +1,6 @@
# Latent Consistency Distillation Example:
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is method to distill latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use the latent consistency distillation to distill stable-diffusion-v1.5 for less timestep inference.
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is a method to distill a latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use latent consistency distillation to distill stable-diffusion-v1.5 for inference with few timesteps.
## Full model distillation
@@ -24,7 +24,7 @@ Then cd in the example folder and run
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
And initialize an [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
@@ -46,12 +46,16 @@ write_basic_config()
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
#### Example with LAION-A6+ dataset
#### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example, and for illustrative purposes only. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/). You may also need to search the hyperparameter space according to the dataset you use.
```bash
runwayml/stable-diffusion-v1-5
PROGRAM="train_lcm_distill_sd_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_sd_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
--resolution=512 \
@@ -59,7 +63,7 @@ PROGRAM="train_lcm_distill_sd_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -69,19 +73,23 @@ PROGRAM="train_lcm_distill_sd_wds.py \
--resume_from_checkpoint=latest \
--report_to=wandb \
--seed=453645634 \
--push_to_hub \
--push_to_hub
```
## LCM-LoRA
Instead of fine-tuning the full model, we can also just train a LoRA that can be injected into any SDXL model.
### Example with LAION-A6+ dataset
### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/).
```bash
runwayml/stable-diffusion-v1-5
PROGRAM="train_lcm_distill_lora_sd_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_lora_sd_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
--resolution=512 \
@@ -90,7 +98,7 @@ PROGRAM="train_lcm_distill_lora_sd_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -1,6 +1,6 @@
# Latent Consistency Distillation Example:
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is method to distill latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use the latent consistency distillation to distill SDXL for less timestep inference.
[Latent Consistency Models (LCMs)](https://arxiv.org/abs/2310.04378) is a method to distill a latent diffusion model to enable swift inference with minimal steps. This example demonstrates how to use latent consistency distillation to distill SDXL for inference with few timesteps.
## Full model distillation
@@ -24,7 +24,7 @@ Then cd in the example folder and run
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
And initialize an [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
@@ -46,12 +46,16 @@ write_basic_config()
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
#### Example with LAION-A6+ dataset
#### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example, and for illustrative purposes only. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/). You may also need to search the hyperparameter space according to the dataset you use.
```bash
export MODEL_DIR="stabilityai/stable-diffusion-xl-base-1.0"
PROGRAM="train_lcm_distill_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_NAME \
--pretrained_vae_model_name_or_path=madebyollin/sdxl-vae-fp16-fix \
--output_dir=$OUTPUT_DIR \
--mixed_precision=fp16 \
@@ -60,7 +64,7 @@ PROGRAM="train_lcm_distill_sdxl_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -77,11 +81,15 @@ PROGRAM="train_lcm_distill_sdxl_wds.py \
Instead of fine-tuning the full model, we can also just train a LoRA that can be injected into any SDXL model.
### Example with LAION-A6+ dataset
### Example
The following uses the [Conceptual Captions 12M (CC12M) dataset](https://github.com/google-research-datasets/conceptual-12m) as an example. For best results you may consider large and high-quality text-image datasets such as [LAION](https://laion.ai/blog/laion-400-open-dataset/).
```bash
export MODEL_DIR="stabilityai/stable-diffusion-xl-base-1.0"
PROGRAM="train_lcm_distill_lora_sdxl_wds.py \
export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export OUTPUT_DIR="path/to/saved/model"
accelerate launch train_lcm_distill_lora_sdxl_wds.py \
--pretrained_teacher_model=$MODEL_DIR \
--pretrained_vae_model_name_or_path=madebyollin/sdxl-vae-fp16-fix \
--output_dir=$OUTPUT_DIR \
@@ -92,7 +100,7 @@ PROGRAM="train_lcm_distill_lora_sdxl_wds.py \
--max_train_steps=1000 \
--max_train_samples=4000000 \
--dataloader_num_workers=8 \
--train_shards_path_or_url='pipe:aws s3 cp s3://muse-datasets/laion-aesthetic6plus-min512-data/{00000..01210}.tar -' \
--train_shards_path_or_url="pipe:curl -L -s https://huggingface.co/datasets/laion/conceptual-captions-12m-webdataset/resolve/main/data/{00000..01099}.tar?download=true" \
--validation_steps=200 \
--checkpointing_steps=200 --checkpoints_total_limit=10 \
--train_batch_size=12 \
@@ -71,7 +71,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.18.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -423,7 +423,7 @@ def import_model_class_from_model_name_or_path(
pretrained_model_name_or_path: str, revision: str, subfolder: str = "text_encoder"
):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path, subfolder=subfolder, revision=revision, use_auth_token=True
pretrained_model_name_or_path, subfolder=subfolder, revision=revision
)
model_class = text_encoder_config.architectures[0]
@@ -1123,7 +1123,7 @@ def main(args):
for epoch in range(first_epoch, args.num_train_epochs):
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
image, text, _, _ = batch
image, text = batch
image = image.to(accelerator.device, non_blocking=True)
encoded_text = compute_embeddings_fn(text)
@@ -68,11 +68,16 @@ from diffusers.utils.import_utils import is_xformers_available
MAX_SEQ_LENGTH = 77
# Adjust for your dataset
WDS_JSON_WIDTH = "width" # original_width for LAION
WDS_JSON_HEIGHT = "height" # original_height for LAION
MIN_SIZE = 700 # ~960 for LAION, ideal: 1024 if the dataset contains large images
if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.18.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -146,10 +151,10 @@ class WebdatasetFilter:
try:
if "json" in x:
x_json = json.loads(x["json"])
filter_size = (x_json.get("original_width", 0.0) or 0.0) >= self.min_size and x_json.get(
"original_height", 0
filter_size = (x_json.get(WDS_JSON_WIDTH, 0.0) or 0.0) >= self.min_size and x_json.get(
WDS_JSON_HEIGHT, 0
) >= self.min_size
filter_watermark = (x_json.get("pwatermark", 1.0) or 1.0) <= self.max_pwatermark
filter_watermark = (x_json.get("pwatermark", 0.0) or 0.0) <= self.max_pwatermark
return filter_size and filter_watermark
else:
return False
@@ -180,7 +185,7 @@ class Text2ImageDataset:
if use_fix_crop_and_size:
return (resolution, resolution)
else:
return (int(json.get("original_width", 0.0)), int(json.get("original_height", 0.0)))
return (int(json.get(WDS_JSON_WIDTH, 0.0)), int(json.get(WDS_JSON_HEIGHT, 0.0)))
def transform(example):
# resize image
@@ -212,7 +217,7 @@ class Text2ImageDataset:
pipeline = [
wds.ResampledShards(train_shards_path_or_url),
tarfile_to_samples_nothrow,
wds.select(WebdatasetFilter(min_size=960)),
wds.select(WebdatasetFilter(min_size=MIN_SIZE)),
wds.shuffle(shuffle_buffer_size),
*processing_pipeline,
wds.batched(per_gpu_batch_size, partial=False, collation_fn=default_collate),
@@ -392,7 +397,7 @@ def import_model_class_from_model_name_or_path(
pretrained_model_name_or_path: str, revision: str, subfolder: str = "text_encoder"
):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path, subfolder=subfolder, revision=revision, use_auth_token=True
pretrained_model_name_or_path, subfolder=subfolder, revision=revision
)
model_class = text_encoder_config.architectures[0]
@@ -70,7 +70,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.18.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -400,7 +400,7 @@ def import_model_class_from_model_name_or_path(
pretrained_model_name_or_path: str, revision: str, subfolder: str = "text_encoder"
):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path, subfolder=subfolder, revision=revision, use_auth_token=True
pretrained_model_name_or_path, subfolder=subfolder, revision=revision
)
model_class = text_encoder_config.architectures[0]
@@ -657,6 +657,15 @@ def parse_args():
default=0.001,
help="The huber loss parameter. Only used if `--loss_type=huber`.",
)
parser.add_argument(
"--unet_time_cond_proj_dim",
type=int,
default=256,
help=(
"The dimension of the guidance scale embedding in the U-Net, which will be used if the teacher U-Net"
" does not have `time_cond_proj_dim` set."
),
)
# ----Exponential Moving Average (EMA)----
parser.add_argument(
"--ema_decay",
@@ -1097,7 +1106,7 @@ def main(args):
for epoch in range(first_epoch, args.num_train_epochs):
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
image, text, _, _ = batch
image, text = batch
image = image.to(accelerator.device, non_blocking=True)
encoded_text = compute_embeddings_fn(text)
@@ -1138,7 +1147,7 @@ def main(args):
# 20.4.6. Sample a random guidance scale w from U[w_min, w_max] and embed it
w = (args.w_max - args.w_min) * torch.rand((bsz,)) + args.w_min
w_embedding = guidance_scale_embedding(w, embedding_dim=args.unet_time_cond_proj_dim)
w_embedding = guidance_scale_embedding(w, embedding_dim=unet.config.time_cond_proj_dim)
w = w.reshape(bsz, 1, 1, 1)
# Move to U-Net device and dtype
w = w.to(device=latents.device, dtype=latents.dtype)
@@ -67,11 +67,16 @@ from diffusers.utils.import_utils import is_xformers_available
MAX_SEQ_LENGTH = 77
# Adjust for your dataset
WDS_JSON_WIDTH = "width" # original_width for LAION
WDS_JSON_HEIGHT = "height" # original_height for LAION
MIN_SIZE = 700 # ~960 for LAION, ideal: 1024 if the dataset contains large images
if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.18.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -128,10 +133,10 @@ class WebdatasetFilter:
try:
if "json" in x:
x_json = json.loads(x["json"])
filter_size = (x_json.get("original_width", 0.0) or 0.0) >= self.min_size and x_json.get(
"original_height", 0
filter_size = (x_json.get(WDS_JSON_WIDTH, 0.0) or 0.0) >= self.min_size and x_json.get(
WDS_JSON_HEIGHT, 0
) >= self.min_size
filter_watermark = (x_json.get("pwatermark", 1.0) or 1.0) <= self.max_pwatermark
filter_watermark = (x_json.get("pwatermark", 0.0) or 0.0) <= self.max_pwatermark
return filter_size and filter_watermark
else:
return False
@@ -162,7 +167,7 @@ class Text2ImageDataset:
if use_fix_crop_and_size:
return (resolution, resolution)
else:
return (int(json.get("original_width", 0.0)), int(json.get("original_height", 0.0)))
return (int(json.get(WDS_JSON_WIDTH, 0.0)), int(json.get(WDS_JSON_HEIGHT, 0.0)))
def transform(example):
# resize image
@@ -194,7 +199,7 @@ class Text2ImageDataset:
pipeline = [
wds.ResampledShards(train_shards_path_or_url),
tarfile_to_samples_nothrow,
wds.select(WebdatasetFilter(min_size=960)),
wds.select(WebdatasetFilter(min_size=MIN_SIZE)),
wds.shuffle(shuffle_buffer_size),
*processing_pipeline,
wds.batched(per_gpu_batch_size, partial=False, collation_fn=default_collate),
@@ -414,7 +419,7 @@ def import_model_class_from_model_name_or_path(
pretrained_model_name_or_path: str, revision: str, subfolder: str = "text_encoder"
):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path, subfolder=subfolder, revision=revision, use_auth_token=True
pretrained_model_name_or_path, subfolder=subfolder, revision=revision
)
model_class = text_encoder_config.architectures[0]
@@ -677,6 +682,15 @@ def parse_args():
default=0.001,
help="The huber loss parameter. Only used if `--loss_type=huber`.",
)
parser.add_argument(
"--unet_time_cond_proj_dim",
type=int,
default=256,
help=(
"The dimension of the guidance scale embedding in the U-Net, which will be used if the teacher U-Net"
" does not have `time_cond_proj_dim` set."
),
)
# ----Exponential Moving Average (EMA)----
parser.add_argument(
"--ema_decay",
@@ -1233,6 +1247,7 @@ def main(args):
# 20.4.6. Sample a random guidance scale w from U[w_min, w_max] and embed it
w = (args.w_max - args.w_min) * torch.rand((bsz,)) + args.w_min
w_embedding = guidance_scale_embedding(w, embedding_dim=unet.config.time_cond_proj_dim)
w = w.reshape(bsz, 1, 1, 1)
w = w.to(device=latents.device, dtype=latents.dtype)
@@ -1243,7 +1258,7 @@ def main(args):
noise_pred = unet(
noisy_model_input,
start_timesteps,
timestep_cond=None,
timestep_cond=w_embedding,
encoder_hidden_states=prompt_embeds.float(),
added_cond_kwargs=encoded_text,
).sample
@@ -1308,7 +1323,7 @@ def main(args):
target_noise_pred = target_unet(
x_prev.float(),
timesteps,
timestep_cond=None,
timestep_cond=w_embedding,
encoder_hidden_states=prompt_embeds.float(),
added_cond_kwargs=encoded_text,
).sample
+120
View File
@@ -0,0 +1,120 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import sys
import tempfile
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class ControlNet(ExamplesTestsAccelerate):
def test_controlnet_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/controlnet/train_controlnet.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/fill10
--output_dir={tmpdir}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
--controlnet_model_name_or_path=hf-internal-testing/tiny-controlnet
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_controlnet_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/controlnet/train_controlnet.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/fill10
--output_dir={tmpdir}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--controlnet_model_name_or_path=hf-internal-testing/tiny-controlnet
--max_train_steps=9
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4", "checkpoint-6", "checkpoint-8"},
)
resume_run_args = f"""
examples/controlnet/train_controlnet.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/fill10
--output_dir={tmpdir}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--controlnet_model_name_or_path=hf-internal-testing/tiny-controlnet
--max_train_steps=11
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-8
--checkpoints_total_limit=3
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-8", "checkpoint-10", "checkpoint-12"},
)
class ControlNetSDXL(ExamplesTestsAccelerate):
def test_controlnet_sdxl(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/controlnet/train_controlnet_sdxl.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-xl-pipe
--dataset_name=hf-internal-testing/fill10
--output_dir={tmpdir}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--controlnet_model_name_or_path=hf-internal-testing/tiny-controlnet-sdxl
--max_train_steps=9
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "diffusion_pytorch_model.safetensors")))
+14 -8
View File
@@ -56,7 +56,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -86,6 +86,7 @@ def log_validation(vae, text_encoder, tokenizer, unet, controlnet, args, acceler
controlnet=controlnet,
safety_checker=None,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline.scheduler = UniPCMultistepScheduler.from_config(pipeline.scheduler.config)
@@ -249,10 +250,13 @@ def parse_args(input_args=None):
type=str,
default=None,
required=False,
help=(
"Revision of pretrained model identifier from huggingface.co/models. Trainable model components should be"
" float32 precision."
),
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--tokenizer_name",
@@ -767,11 +771,13 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
if args.controlnet_model_name_or_path:
+1 -1
View File
@@ -59,7 +59,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = logging.getLogger(__name__)
+21 -10
View File
@@ -58,7 +58,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -74,6 +74,7 @@ def log_validation(vae, unet, controlnet, args, accelerator, weight_dtype, step)
unet=unet,
controlnet=controlnet,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline.scheduler = UniPCMultistepScheduler.from_config(pipeline.scheduler.config)
@@ -243,15 +244,18 @@ def parse_args(input_args=None):
help="Path to pretrained controlnet model or model identifier from huggingface.co/models."
" If not specified controlnet weights are initialized from unet.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--revision",
type=str,
default=None,
required=False,
help=(
"Revision of pretrained model identifier from huggingface.co/models. Trainable model components should be"
" float32 precision."
),
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--tokenizer_name",
@@ -793,10 +797,16 @@ def main(args):
# Load the tokenizers
tokenizer_one = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer",
revision=args.revision,
use_fast=False,
)
tokenizer_two = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer_2", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer_2",
revision=args.revision,
use_fast=False,
)
# import correct text encoder classes
@@ -810,10 +820,10 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder_one = text_encoder_cls_one.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
text_encoder_two = text_encoder_cls_two.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision, variant=args.variant
)
vae_path = (
args.pretrained_model_name_or_path
@@ -824,9 +834,10 @@ def main(args):
vae_path,
subfolder="vae" if args.pretrained_vae_model_name_or_path is None else None,
revision=args.revision,
variant=args.variant,
)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
if args.controlnet_model_name_or_path:
@@ -0,0 +1,130 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import sys
import tempfile
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class CustomDiffusion(ExamplesTestsAccelerate):
def test_custom_diffusion(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/custom_diffusion/train_custom_diffusion.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt <new1>
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 1.0e-05
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--modifier_token <new1>
--no_safe_serialization
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_custom_diffusion_weights.bin")))
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "<new1>.bin")))
def test_custom_diffusion_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/custom_diffusion/train_custom_diffusion.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=<new1>
--resolution=64
--train_batch_size=1
--modifier_token=<new1>
--dataloader_num_workers=0
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
--no_safe_serialization
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_custom_diffusion_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/custom_diffusion/train_custom_diffusion.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=<new1>
--resolution=64
--train_batch_size=1
--modifier_token=<new1>
--dataloader_num_workers=0
--max_train_steps=9
--checkpointing_steps=2
--no_safe_serialization
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4", "checkpoint-6", "checkpoint-8"},
)
resume_run_args = f"""
examples/custom_diffusion/train_custom_diffusion.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=<new1>
--resolution=64
--train_batch_size=1
--modifier_token=<new1>
--dataloader_num_workers=0
--max_train_steps=11
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-8
--checkpoints_total_limit=3
--no_safe_serialization
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-6", "checkpoint-8", "checkpoint-10"},
)
@@ -62,7 +62,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -332,6 +332,12 @@ def parse_args(input_args=None):
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--tokenizer_name",
type=str,
@@ -740,6 +746,7 @@ def main(args):
torch_dtype=torch_dtype,
safety_checker=None,
revision=args.revision,
variant=args.variant,
)
pipeline.set_progress_bar_config(disable=True)
@@ -801,11 +808,13 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
# Adding a modifier token which is optimized ####
@@ -1229,6 +1238,7 @@ def main(args):
text_encoder=accelerator.unwrap_model(text_encoder),
tokenizer=tokenizer,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
@@ -1278,7 +1288,7 @@ def main(args):
# Final inference
# Load previous pipeline
pipeline = DiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path, revision=args.revision, torch_dtype=weight_dtype
args.pretrained_model_name_or_path, revision=args.revision, variant=args.variant, torch_dtype=weight_dtype
)
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
pipeline = pipeline.to(accelerator.device)
+1
View File
@@ -44,6 +44,7 @@ write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Dog toy example
+1
View File
@@ -47,6 +47,7 @@ write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Dog toy example
+1
View File
@@ -4,3 +4,4 @@ transformers>=4.25.1
ftfy
tensorboard
Jinja2
peft==0.7.0
@@ -4,3 +4,4 @@ transformers>=4.25.1
ftfy
tensorboard
Jinja2
peft==0.7.0
+230
View File
@@ -0,0 +1,230 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import shutil
import sys
import tempfile
from diffusers import DiffusionPipeline, UNet2DConditionModel
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class DreamBooth(ExamplesTestsAccelerate):
def test_dreambooth(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "unet", "diffusion_pytorch_model.safetensors")))
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "scheduler", "scheduler_config.json")))
def test_dreambooth_if(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path hf-internal-testing/tiny-if-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--pre_compute_text_embeddings
--tokenizer_max_length=77
--text_encoder_use_attention_mask
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "unet", "diffusion_pytorch_model.safetensors")))
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "scheduler", "scheduler_config.json")))
def test_dreambooth_checkpointing(self):
instance_prompt = "photo"
pretrained_model_name_or_path = "hf-internal-testing/tiny-stable-diffusion-pipe"
with tempfile.TemporaryDirectory() as tmpdir:
# Run training script with checkpointing
# max_train_steps == 5, checkpointing_steps == 2
# Should create checkpoints at steps 2, 4
initial_run_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path {pretrained_model_name_or_path}
--instance_data_dir docs/source/en/imgs
--instance_prompt {instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 5
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--checkpointing_steps=2
--seed=0
""".split()
run_command(self._launch_args + initial_run_args)
# check can run the original fully trained output pipeline
pipe = DiffusionPipeline.from_pretrained(tmpdir, safety_checker=None)
pipe(instance_prompt, num_inference_steps=2)
# check checkpoint directories exist
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
# check can run an intermediate checkpoint
unet = UNet2DConditionModel.from_pretrained(tmpdir, subfolder="checkpoint-2/unet")
pipe = DiffusionPipeline.from_pretrained(pretrained_model_name_or_path, unet=unet, safety_checker=None)
pipe(instance_prompt, num_inference_steps=2)
# Remove checkpoint 2 so that we can check only later checkpoints exist after resuming
shutil.rmtree(os.path.join(tmpdir, "checkpoint-2"))
# Run training script for 7 total steps resuming from checkpoint 4
resume_run_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path {pretrained_model_name_or_path}
--instance_data_dir docs/source/en/imgs
--instance_prompt {instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 7
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-4
--seed=0
""".split()
run_command(self._launch_args + resume_run_args)
# check can run new fully trained pipeline
pipe = DiffusionPipeline.from_pretrained(tmpdir, safety_checker=None)
pipe(instance_prompt, num_inference_steps=2)
# check old checkpoints do not exist
self.assertFalse(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
# check new checkpoints exist
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-6")))
def test_dreambooth_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_dreambooth_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=9
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4", "checkpoint-6", "checkpoint-8"},
)
resume_run_args = f"""
examples/dreambooth/train_dreambooth.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=11
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-8
--checkpoints_total_limit=3
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-6", "checkpoint-8", "checkpoint-10"},
)
+388
View File
@@ -0,0 +1,388 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import sys
import tempfile
import safetensors
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
from diffusers import DiffusionPipeline # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class DreamBoothLoRA(ExamplesTestsAccelerate):
def test_dreambooth_lora(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
# when not training the text encoder, all the parameters in the state dict should start
# with `"unet"` in their names.
starts_with_unet = all(key.startswith("unet") for key in lora_state_dict.keys())
self.assertTrue(starts_with_unet)
def test_dreambooth_lora_with_text_encoder(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--train_text_encoder
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# check `text_encoder` is present at all.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
keys = lora_state_dict.keys()
is_text_encoder_present = any(k.startswith("text_encoder") for k in keys)
self.assertTrue(is_text_encoder_present)
# the names of the keys of the state dict should either start with `unet`
# or `text_encoder`.
is_correct_naming = all(k.startswith("unet") or k.startswith("text_encoder") for k in keys)
self.assertTrue(is_correct_naming)
def test_dreambooth_lora_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_dreambooth_lora_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=9
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4", "checkpoint-6", "checkpoint-8"},
)
resume_run_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--instance_data_dir=docs/source/en/imgs
--output_dir={tmpdir}
--instance_prompt=prompt
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=11
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-8
--checkpoints_total_limit=3
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-6", "checkpoint-8", "checkpoint-10"},
)
def test_dreambooth_lora_if_model(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora.py
--pretrained_model_name_or_path hf-internal-testing/tiny-if-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--pre_compute_text_embeddings
--tokenizer_max_length=77
--text_encoder_use_attention_mask
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
# when not training the text encoder, all the parameters in the state dict should start
# with `"unet"` in their names.
starts_with_unet = all(key.startswith("unet") for key in lora_state_dict.keys())
self.assertTrue(starts_with_unet)
class DreamBoothLoRASDXL(ExamplesTestsAccelerate):
def test_dreambooth_lora_sdxl(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-xl-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
# when not training the text encoder, all the parameters in the state dict should start
# with `"unet"` in their names.
starts_with_unet = all(key.startswith("unet") for key in lora_state_dict.keys())
self.assertTrue(starts_with_unet)
def test_dreambooth_lora_sdxl_with_text_encoder(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-xl-pipe
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--train_text_encoder
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
# when not training the text encoder, all the parameters in the state dict should start
# with `"unet"` or `"text_encoder"` or `"text_encoder_2"` in their names.
keys = lora_state_dict.keys()
starts_with_unet = all(
k.startswith("unet") or k.startswith("text_encoder") or k.startswith("text_encoder_2") for k in keys
)
self.assertTrue(starts_with_unet)
def test_dreambooth_lora_sdxl_custom_captions(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-xl-pipe
--dataset_name hf-internal-testing/dummy_image_text_data
--caption_column text
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
def test_dreambooth_lora_sdxl_text_encoder_custom_captions(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path hf-internal-testing/tiny-stable-diffusion-xl-pipe
--dataset_name hf-internal-testing/dummy_image_text_data
--caption_column text
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--train_text_encoder
""".split()
run_command(self._launch_args + test_args)
def test_dreambooth_lora_sdxl_checkpointing_checkpoints_total_limit(self):
pipeline_path = "hf-internal-testing/tiny-stable-diffusion-xl-pipe"
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path {pipeline_path}
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 7
--checkpointing_steps=2
--checkpoints_total_limit=2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
pipe = DiffusionPipeline.from_pretrained(pipeline_path)
pipe.load_lora_weights(tmpdir)
pipe("a prompt", num_inference_steps=2)
# check checkpoint directories exist
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
# checkpoint-2 should have been deleted
{"checkpoint-4", "checkpoint-6"},
)
def test_dreambooth_lora_sdxl_text_encoder_checkpointing_checkpoints_total_limit(self):
pipeline_path = "hf-internal-testing/tiny-stable-diffusion-xl-pipe"
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/dreambooth/train_dreambooth_lora_sdxl.py
--pretrained_model_name_or_path {pipeline_path}
--instance_data_dir docs/source/en/imgs
--instance_prompt photo
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 7
--checkpointing_steps=2
--checkpoints_total_limit=2
--train_text_encoder
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
pipe = DiffusionPipeline.from_pretrained(pipeline_path)
pipe.load_lora_weights(tmpdir)
pipe("a prompt", num_inference_steps=2)
# check checkpoint directories exist
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
# checkpoint-2 should have been deleted
{"checkpoint-4", "checkpoint-6"},
)
+15 -9
View File
@@ -61,7 +61,7 @@ if is_wandb_available():
import wandb
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
@@ -139,6 +139,7 @@ def log_validation(
text_encoder=text_encoder,
unet=accelerator.unwrap_model(unet),
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
**pipeline_args,
)
@@ -239,10 +240,13 @@ def parse_args(input_args=None):
type=str,
default=None,
required=False,
help=(
"Revision of pretrained model identifier from huggingface.co/models. Trainable model components should be"
" float32 precision."
),
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--tokenizer_name",
@@ -296,7 +300,7 @@ def parse_args(input_args=None):
parser.add_argument(
"--output_dir",
type=str,
default="text-inversion-model",
default="dreambooth-model",
help="The output directory where the model predictions and checkpoints will be written.",
)
parser.add_argument("--seed", type=int, default=None, help="A seed for reproducible training.")
@@ -859,6 +863,7 @@ def main(args):
torch_dtype=torch_dtype,
safety_checker=None,
revision=args.revision,
variant=args.variant,
)
pipeline.set_progress_bar_config(disable=True)
@@ -912,18 +917,18 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
if model_has_vae(args):
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision, variant=args.variant
)
else:
vae = None
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
@@ -1379,6 +1384,7 @@ def main(args):
args.pretrained_model_name_or_path,
unet=accelerator.unwrap_model(unet),
revision=args.revision,
variant=args.variant,
**pipeline_args,
)
+9 -3
View File
@@ -35,7 +35,7 @@ from diffusers.utils import check_min_version
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
# Cache compiled models across invocations of this script.
cc.initialize_cache(os.path.expanduser("~/.cache/jax/compilation_cache"))
@@ -460,7 +460,10 @@ def main():
# Load models and create wrapper for stable diffusion
text_encoder = FlaxCLIPTextModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", dtype=weight_dtype, revision=args.revision
args.pretrained_model_name_or_path,
subfolder="text_encoder",
dtype=weight_dtype,
revision=args.revision,
)
vae, vae_params = FlaxAutoencoderKL.from_pretrained(
vae_arg,
@@ -468,7 +471,10 @@ def main():
**vae_kwargs,
)
unet, unet_params = FlaxUNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", dtype=weight_dtype, revision=args.revision
args.pretrained_model_name_or_path,
subfolder="unet",
dtype=weight_dtype,
revision=args.revision,
)
# Optimization
+73 -104
View File
@@ -16,7 +16,6 @@
import argparse
import copy
import gc
import itertools
import logging
import math
import os
@@ -35,6 +34,8 @@ from accelerate.utils import ProjectConfiguration, set_seed
from huggingface_hub import create_repo, upload_folder
from huggingface_hub.utils import insecure_hashlib
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from PIL import Image
from PIL.ImageOps import exif_transpose
from torch.utils.data import Dataset
@@ -52,24 +53,50 @@ from diffusers import (
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.models.attention_processor import (
AttnAddedKVProcessor,
AttnAddedKVProcessor2_0,
SlicedAttnAddedKVProcessor,
)
from diffusers.models.lora import LoRALinearLayer, text_encoder_lora_state_dict
from diffusers.optimization import get_scheduler
from diffusers.training_utils import unet_lora_state_dict
from diffusers.utils import check_min_version, is_wandb_available
from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
# TODO: This function should be removed once training scripts are rewritten in PEFT
def text_encoder_lora_state_dict(text_encoder):
state_dict = {}
def text_encoder_attn_modules(text_encoder):
from transformers import CLIPTextModel, CLIPTextModelWithProjection
attn_modules = []
if isinstance(text_encoder, (CLIPTextModel, CLIPTextModelWithProjection)):
for i, layer in enumerate(text_encoder.text_model.encoder.layers):
name = f"text_model.encoder.layers.{i}.self_attn"
mod = layer.self_attn
attn_modules.append((name, mod))
return attn_modules
for name, module in text_encoder_attn_modules(text_encoder):
for k, v in module.q_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.q_proj.lora_linear_layer.{k}"] = v
for k, v in module.k_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.k_proj.lora_linear_layer.{k}"] = v
for k, v in module.v_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.v_proj.lora_linear_layer.{k}"] = v
for k, v in module.out_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.out_proj.lora_linear_layer.{k}"] = v
return state_dict
def save_model_card(
repo_id: str,
images=None,
@@ -150,6 +177,12 @@ def parse_args(input_args=None):
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--tokenizer_name",
type=str,
@@ -717,6 +750,7 @@ def main(args):
torch_dtype=torch_dtype,
safety_checker=None,
revision=args.revision,
variant=args.variant,
)
pipeline.set_progress_bar_config(disable=True)
@@ -770,11 +804,11 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
try:
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision, variant=args.variant
)
except OSError:
# IF does not have a VAE so let's just set it to None
@@ -782,7 +816,7 @@ def main(args):
vae = None
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
# We only train the additional adapter LoRA layers
@@ -824,79 +858,19 @@ def main(args):
text_encoder.gradient_checkpointing_enable()
# now we will add new LoRA weights to the attention layers
# It's important to realize here how many attention weights will be added and of which sizes
# The sizes of the attention layers consist only of two different variables:
# 1) - the "hidden_size", which is increased according to `unet.config.block_out_channels`.
# 2) - the "cross attention size", which is set to `unet.config.cross_attention_dim`.
unet_lora_config = LoraConfig(
r=args.rank,
init_lora_weights="gaussian",
target_modules=["to_k", "to_q", "to_v", "to_out.0", "add_k_proj", "add_v_proj"],
)
unet.add_adapter(unet_lora_config)
# Let's first see how many attention processors we will have to set.
# For Stable Diffusion, it should be equal to:
# - down blocks (2x attention layers) * (2x transformer layers) * (3x down blocks) = 12
# - mid blocks (2x attention layers) * (1x transformer layers) * (1x mid blocks) = 2
# - up blocks (2x attention layers) * (3x transformer layers) * (3x up blocks) = 18
# => 32 layers
# Set correct lora layers
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
if isinstance(attn_processor, (AttnAddedKVProcessor, SlicedAttnAddedKVProcessor, AttnAddedKVProcessor2_0)):
attn_module.add_k_proj.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.add_k_proj.in_features,
out_features=attn_module.add_k_proj.out_features,
rank=args.rank,
)
)
attn_module.add_v_proj.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.add_v_proj.in_features,
out_features=attn_module.add_v_proj.out_features,
rank=args.rank,
)
)
unet_lora_parameters.extend(attn_module.add_k_proj.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.add_v_proj.lora_layer.parameters())
# The text encoder comes from 🤗 transformers, so we cannot directly modify it.
# So, instead, we monkey-patch the forward calls of its attention-blocks.
# The text encoder comes from 🤗 transformers, we will also attach adapters to it.
if args.train_text_encoder:
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
text_lora_parameters = LoraLoaderMixin._modify_text_encoder(text_encoder, dtype=torch.float32, rank=args.rank)
text_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["q_proj", "k_proj", "v_proj", "out_proj"]
)
text_encoder.add_adapter(text_lora_config)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
def save_model_hook(models, weights, output_dir):
@@ -908,9 +882,9 @@ def main(args):
for model in models:
if isinstance(model, type(accelerator.unwrap_model(unet))):
unet_lora_layers_to_save = unet_lora_state_dict(model)
unet_lora_layers_to_save = get_peft_model_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder))):
text_encoder_lora_layers_to_save = text_encoder_lora_state_dict(model)
text_encoder_lora_layers_to_save = get_peft_model_state_dict(model)
else:
raise ValueError(f"unexpected save model: {model.__class__}")
@@ -970,11 +944,10 @@ def main(args):
optimizer_class = torch.optim.AdamW
# Optimizer creation
params_to_optimize = (
itertools.chain(unet_lora_parameters, text_lora_parameters)
if args.train_text_encoder
else unet_lora_parameters
)
params_to_optimize = list(filter(lambda p: p.requires_grad, unet.parameters()))
if args.train_text_encoder:
params_to_optimize = params_to_optimize + list(filter(lambda p: p.requires_grad, text_encoder.parameters()))
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
@@ -1217,12 +1190,7 @@ def main(args):
accelerator.backward(loss)
if accelerator.sync_gradients:
params_to_clip = (
itertools.chain(unet_lora_parameters, text_lora_parameters)
if args.train_text_encoder
else unet_lora_parameters
)
accelerator.clip_grad_norm_(params_to_clip, args.max_grad_norm)
accelerator.clip_grad_norm_(params_to_optimize, args.max_grad_norm)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
@@ -1277,6 +1245,7 @@ def main(args):
unet=accelerator.unwrap_model(unet),
text_encoder=None if args.pre_compute_text_embeddings else accelerator.unwrap_model(text_encoder),
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
@@ -1344,25 +1313,25 @@ def main(args):
if accelerator.is_main_process:
unet = accelerator.unwrap_model(unet)
unet = unet.to(torch.float32)
unet_lora_layers = unet_lora_state_dict(unet)
if text_encoder is not None and args.train_text_encoder:
unet_lora_state_dict = get_peft_model_state_dict(unet)
if args.train_text_encoder:
text_encoder = accelerator.unwrap_model(text_encoder)
text_encoder = text_encoder.to(torch.float32)
text_encoder_lora_layers = text_encoder_lora_state_dict(text_encoder)
text_encoder_state_dict = get_peft_model_state_dict(text_encoder)
else:
text_encoder_lora_layers = None
text_encoder_state_dict = None
LoraLoaderMixin.save_lora_weights(
save_directory=args.output_dir,
unet_lora_layers=unet_lora_layers,
text_encoder_lora_layers=text_encoder_lora_layers,
unet_lora_layers=unet_lora_state_dict,
text_encoder_lora_layers=text_encoder_state_dict,
)
# Final inference
# Load previous pipeline
pipeline = DiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path, revision=args.revision, torch_dtype=weight_dtype
args.pretrained_model_name_or_path, revision=args.revision, variant=args.variant, torch_dtype=weight_dtype
)
# We train on the simplified learning objective. If we were previously predicting a variance, we need the scheduler to ignore it
+503 -159
View File
@@ -34,6 +34,8 @@ from accelerate.utils import DistributedDataParallelKwargs, ProjectConfiguration
from huggingface_hub import create_repo, upload_folder
from huggingface_hub.utils import insecure_hashlib
from packaging import version
from peft import LoraConfig
from peft.utils import get_peft_model_state_dict
from PIL import Image
from PIL.ImageOps import exif_transpose
from torch.utils.data import Dataset
@@ -50,50 +52,112 @@ from diffusers import (
UNet2DConditionModel,
)
from diffusers.loaders import LoraLoaderMixin
from diffusers.models.lora import LoRALinearLayer, text_encoder_lora_state_dict
from diffusers.optimization import get_scheduler
from diffusers.training_utils import unet_lora_state_dict
from diffusers.training_utils import compute_snr
from diffusers.utils import check_min_version, is_wandb_available
from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__)
# TODO: This function should be removed once training scripts are rewritten in PEFT
def text_encoder_lora_state_dict(text_encoder):
state_dict = {}
def text_encoder_attn_modules(text_encoder):
from transformers import CLIPTextModel, CLIPTextModelWithProjection
attn_modules = []
if isinstance(text_encoder, (CLIPTextModel, CLIPTextModelWithProjection)):
for i, layer in enumerate(text_encoder.text_model.encoder.layers):
name = f"text_model.encoder.layers.{i}.self_attn"
mod = layer.self_attn
attn_modules.append((name, mod))
return attn_modules
for name, module in text_encoder_attn_modules(text_encoder):
for k, v in module.q_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.q_proj.lora_linear_layer.{k}"] = v
for k, v in module.k_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.k_proj.lora_linear_layer.{k}"] = v
for k, v in module.v_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.v_proj.lora_linear_layer.{k}"] = v
for k, v in module.out_proj.lora_linear_layer.state_dict().items():
state_dict[f"{name}.out_proj.lora_linear_layer.{k}"] = v
return state_dict
def save_model_card(
repo_id: str, images=None, base_model=str, train_text_encoder=False, prompt=str, repo_folder=None, vae_path=None
repo_id: str,
images=None,
base_model=str,
train_text_encoder=False,
instance_prompt=str,
validation_prompt=str,
repo_folder=None,
vae_path=None,
):
img_str = ""
img_str = "widget:\n" if images else ""
for i, image in enumerate(images):
image.save(os.path.join(repo_folder, f"image_{i}.png"))
img_str += f"![img_{i}](./image_{i}.png)\n"
img_str += f"""
- text: '{validation_prompt if validation_prompt else ' ' }'
output:
url:
"image_{i}.png"
"""
yaml = f"""
---
license: openrail++
base_model: {base_model}
instance_prompt: {prompt}
tags:
- stable-diffusion-xl
- stable-diffusion-xl-diffusers
- text-to-image
- diffusers
- lora
inference: true
- template:sd-lora
{img_str}
base_model: {base_model}
instance_prompt: {instance_prompt}
license: openrail++
---
"""
model_card = f"""
# LoRA DreamBooth - {repo_id}
These are LoRA adaption weights for {base_model}. The weights were trained on {prompt} using [DreamBooth](https://dreambooth.github.io/). You can find some example images in the following. \n
{img_str}
model_card = f"""
# SDXL LoRA DreamBooth - {repo_id}
<Gallery />
## Model description
These are {repo_id} LoRA adaption weights for {base_model}.
The weights were trained using [DreamBooth](https://dreambooth.github.io/).
LoRA for the text encoder was enabled: {train_text_encoder}.
Special VAE used for training: {vae_path}.
## Trigger words
You should use {instance_prompt} to trigger the image generation.
## Download model
Weights for this model are available in Safetensors format.
[Download]({repo_id}/tree/main) them in the Files & versions tab.
"""
with open(os.path.join(repo_folder, "README.md"), "w") as f:
f.write(yaml + model_card)
@@ -141,13 +205,59 @@ def parse_args(input_args=None):
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--dataset_name",
type=str,
default=None,
help=(
"The name of the Dataset (from the HuggingFace hub) containing the training data of instance images (could be your own, possibly private,"
" dataset). It can also be a path pointing to a local copy of a dataset in your filesystem,"
" or to a folder containing files that 🤗 Datasets can understand."
),
)
parser.add_argument(
"--dataset_config_name",
type=str,
default=None,
help="The config of the Dataset, leave as None if there's only one config.",
)
parser.add_argument(
"--instance_data_dir",
type=str,
default=None,
required=True,
help="A folder containing the training data of instance images.",
help=("A folder containing the training data. "),
)
parser.add_argument(
"--cache_dir",
type=str,
default=None,
help="The directory where the downloaded models and datasets will be stored.",
)
parser.add_argument(
"--image_column",
type=str,
default="image",
help="The column of the dataset containing the target image. By "
"default, the standard Image Dataset maps out 'file_name' "
"to 'image'.",
)
parser.add_argument(
"--caption_column",
type=str,
default=None,
help="The column of the dataset containing the instance prompt for each image",
)
parser.add_argument("--repeats", type=int, default=1, help="How many times to repeat the training data.")
parser.add_argument(
"--class_data_dir",
type=str,
@@ -160,7 +270,7 @@ def parse_args(input_args=None):
type=str,
default=None,
required=True,
help="The prompt with identifier specifying the instance",
help="The prompt with identifier specifying the instance, e.g. 'photo of a TOK dog', 'in the style of TOK'",
)
parser.add_argument(
"--class_prompt",
@@ -299,9 +409,16 @@ def parse_args(input_args=None):
parser.add_argument(
"--learning_rate",
type=float,
default=5e-4,
default=1e-4,
help="Initial learning rate (after the potential warmup period) to use.",
)
parser.add_argument(
"--text_encoder_lr",
type=float,
default=5e-6,
help="Text encoder learning rate to use.",
)
parser.add_argument(
"--scale_lr",
action="store_true",
@@ -317,6 +434,14 @@ def parse_args(input_args=None):
' "constant", "constant_with_warmup"]'
),
)
parser.add_argument(
"--snr_gamma",
type=float,
default=None,
help="SNR weighting gamma to be used if rebalancing the loss. Recommended value is 5.0. "
"More details here: https://arxiv.org/abs/2303.09556.",
)
parser.add_argument(
"--lr_warmup_steps", type=int, default=500, help="Number of steps for the warmup in the lr scheduler."
)
@@ -335,13 +460,59 @@ def parse_args(input_args=None):
"Number of subprocesses to use for data loading. 0 means that the data will be loaded in the main process."
),
)
parser.add_argument(
"--use_8bit_adam", action="store_true", help="Whether or not to use 8-bit Adam from bitsandbytes."
"--optimizer",
type=str,
default="AdamW",
help=('The optimizer type to use. Choose between ["AdamW", "prodigy"]'),
)
parser.add_argument(
"--use_8bit_adam",
action="store_true",
help="Whether or not to use 8-bit Adam from bitsandbytes. Ignored if optimizer is not set to AdamW",
)
parser.add_argument(
"--adam_beta1", type=float, default=0.9, help="The beta1 parameter for the Adam and Prodigy optimizers."
)
parser.add_argument(
"--adam_beta2", type=float, default=0.999, help="The beta2 parameter for the Adam and Prodigy optimizers."
)
parser.add_argument(
"--prodigy_beta3",
type=float,
default=None,
help="coefficients for computing the Prodidy stepsize using running averages. If set to None, "
"uses the value of square root of beta2. Ignored if optimizer is adamW",
)
parser.add_argument("--prodigy_decouple", type=bool, default=True, help="Use AdamW style decoupled weight decay")
parser.add_argument("--adam_weight_decay", type=float, default=1e-04, help="Weight decay to use for unet params")
parser.add_argument(
"--adam_weight_decay_text_encoder", type=float, default=1e-03, help="Weight decay to use for text_encoder"
)
parser.add_argument(
"--adam_epsilon",
type=float,
default=1e-08,
help="Epsilon value for the Adam optimizer and Prodigy optimizers.",
)
parser.add_argument(
"--prodigy_use_bias_correction",
type=bool,
default=True,
help="Turn on Adam's bias correction. True by default. Ignored if optimizer is adamW",
)
parser.add_argument(
"--prodigy_safeguard_warmup",
type=bool,
default=True,
help="Remove lr from the denominator of D estimate to avoid issues during warm-up stage. True by default. "
"Ignored if optimizer is adamW",
)
parser.add_argument("--adam_beta1", type=float, default=0.9, help="The beta1 parameter for the Adam optimizer.")
parser.add_argument("--adam_beta2", type=float, default=0.999, help="The beta2 parameter for the Adam optimizer.")
parser.add_argument("--adam_weight_decay", type=float, default=1e-2, help="Weight decay to use.")
parser.add_argument("--adam_epsilon", type=float, default=1e-08, help="Epsilon value for the Adam optimizer")
parser.add_argument("--max_grad_norm", default=1.0, type=float, help="Max gradient norm.")
parser.add_argument("--push_to_hub", action="store_true", help="Whether or not to push the model to the Hub.")
parser.add_argument("--hub_token", type=str, default=None, help="The token to use to push to the Model Hub.")
@@ -414,6 +585,12 @@ def parse_args(input_args=None):
else:
args = parser.parse_args()
if args.dataset_name is None and args.instance_data_dir is None:
raise ValueError("Specify either `--dataset_name` or `--instance_data_dir`")
if args.dataset_name is not None and args.instance_data_dir is not None:
raise ValueError("Specify only one of `--dataset_name` or `--instance_data_dir`")
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
if env_local_rank != -1 and env_local_rank != args.local_rank:
args.local_rank = env_local_rank
@@ -442,20 +619,84 @@ class DreamBoothDataset(Dataset):
def __init__(
self,
instance_data_root,
instance_prompt,
class_prompt,
class_data_root=None,
class_num=None,
size=1024,
repeats=1,
center_crop=False,
):
self.size = size
self.center_crop = center_crop
self.instance_data_root = Path(instance_data_root)
if not self.instance_data_root.exists():
raise ValueError("Instance images root doesn't exists.")
self.instance_prompt = instance_prompt
self.custom_instance_prompts = None
self.class_prompt = class_prompt
self.instance_images_path = list(Path(instance_data_root).iterdir())
self.num_instance_images = len(self.instance_images_path)
# if --dataset_name is provided or a metadata jsonl file is provided in the local --instance_data directory,
# we load the training data using load_dataset
if args.dataset_name is not None:
try:
from datasets import load_dataset
except ImportError:
raise ImportError(
"You are trying to load your data using the datasets library. If you wish to train using custom "
"captions please install the datasets library: `pip install datasets`. If you wish to load a "
"local folder containing images only, specify --instance_data_dir instead."
)
# Downloading and loading a dataset from the hub.
# See more about loading custom images at
# https://huggingface.co/docs/datasets/v2.0.0/en/dataset_script
dataset = load_dataset(
args.dataset_name,
args.dataset_config_name,
cache_dir=args.cache_dir,
)
# Preprocessing the datasets.
column_names = dataset["train"].column_names
# 6. Get the column names for input/target.
if args.image_column is None:
image_column = column_names[0]
logger.info(f"image column defaulting to {image_column}")
else:
image_column = args.image_column
if image_column not in column_names:
raise ValueError(
f"`--image_column` value '{args.image_column}' not found in dataset columns. Dataset columns are: {', '.join(column_names)}"
)
instance_images = dataset["train"][image_column]
if args.caption_column is None:
logger.info(
"No caption column provided, defaulting to instance_prompt for all images. If your dataset "
"contains captions/prompts for the images, make sure to specify the "
"column as --caption_column"
)
self.custom_instance_prompts = None
else:
if args.caption_column not in column_names:
raise ValueError(
f"`--caption_column` value '{args.caption_column}' not found in dataset columns. Dataset columns are: {', '.join(column_names)}"
)
custom_instance_prompts = dataset["train"][args.caption_column]
# create final list of captions according to --repeats
self.custom_instance_prompts = []
for caption in custom_instance_prompts:
self.custom_instance_prompts.extend(itertools.repeat(caption, repeats))
else:
self.instance_data_root = Path(instance_data_root)
if not self.instance_data_root.exists():
raise ValueError("Instance images root doesn't exists.")
instance_images = [Image.open(path) for path in list(Path(instance_data_root).iterdir())]
self.custom_instance_prompts = None
self.instance_images = []
for img in instance_images:
self.instance_images.extend(itertools.repeat(img, repeats))
self.num_instance_images = len(self.instance_images)
self._length = self.num_instance_images
if class_data_root is not None:
@@ -484,13 +725,23 @@ class DreamBoothDataset(Dataset):
def __getitem__(self, index):
example = {}
instance_image = Image.open(self.instance_images_path[index % self.num_instance_images])
instance_image = self.instance_images[index % self.num_instance_images]
instance_image = exif_transpose(instance_image)
if not instance_image.mode == "RGB":
instance_image = instance_image.convert("RGB")
example["instance_images"] = self.image_transforms(instance_image)
if self.custom_instance_prompts:
caption = self.custom_instance_prompts[index % self.num_instance_images]
if caption:
example["instance_prompt"] = caption
else:
example["instance_prompt"] = self.instance_prompt
else: # costum prompts were provided, but length does not match size of image dataset
example["instance_prompt"] = self.instance_prompt
if self.class_data_root:
class_image = Image.open(self.class_images_path[index % self.num_class_images])
class_image = exif_transpose(class_image)
@@ -498,22 +749,25 @@ class DreamBoothDataset(Dataset):
if not class_image.mode == "RGB":
class_image = class_image.convert("RGB")
example["class_images"] = self.image_transforms(class_image)
example["class_prompt"] = self.class_prompt
return example
def collate_fn(examples, with_prior_preservation=False):
pixel_values = [example["instance_images"] for example in examples]
prompts = [example["instance_prompt"] for example in examples]
# Concat class and instance examples for prior preservation.
# We do this to avoid doing two forward passes.
if with_prior_preservation:
pixel_values += [example["class_images"] for example in examples]
prompts += [example["class_prompt"] for example in examples]
pixel_values = torch.stack(pixel_values)
pixel_values = pixel_values.to(memory_format=torch.contiguous_format).float()
batch = {"pixel_values": pixel_values}
batch = {"pixel_values": pixel_values, "prompts": prompts}
return batch
@@ -630,6 +884,7 @@ def main(args):
args.pretrained_model_name_or_path,
torch_dtype=torch_dtype,
revision=args.revision,
variant=args.variant,
)
pipeline.set_progress_bar_config(disable=True)
@@ -668,10 +923,16 @@ def main(args):
# Load the tokenizers
tokenizer_one = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer",
revision=args.revision,
use_fast=False,
)
tokenizer_two = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer_2", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer_2",
revision=args.revision,
use_fast=False,
)
# import correct text encoder classes
@@ -685,10 +946,10 @@ def main(args):
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder_one = text_encoder_cls_one.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
text_encoder_two = text_encoder_cls_two.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision, variant=args.variant
)
vae_path = (
args.pretrained_model_name_or_path
@@ -696,10 +957,13 @@ def main(args):
else args.pretrained_vae_model_name_or_path
)
vae = AutoencoderKL.from_pretrained(
vae_path, subfolder="vae" if args.pretrained_vae_model_name_or_path is None else None, revision=args.revision
vae_path,
subfolder="vae" if args.pretrained_vae_model_name_or_path is None else None,
revision=args.revision,
variant=args.variant,
)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
# We only train the additional adapter LoRA layers
@@ -732,7 +996,8 @@ def main(args):
xformers_version = version.parse(xformers.__version__)
if xformers_version == version.parse("0.0.16"):
logger.warn(
"xFormers 0.0.16 cannot be used for training in some GPUs. If you observe problems during training, please update xFormers to at least 0.0.17. See https://huggingface.co/docs/diffusers/main/en/optimization/xformers for more details."
"xFormers 0.0.16 cannot be used for training in some GPUs. If you observe problems during training, "
"please update xFormers to at least 0.0.17. See https://huggingface.co/docs/diffusers/main/en/optimization/xformers for more details."
)
unet.enable_xformers_memory_efficient_attention()
else:
@@ -745,54 +1010,19 @@ def main(args):
text_encoder_two.gradient_checkpointing_enable()
# now we will add new LoRA weights to the attention layers
# Set correct lora layers
unet_lora_parameters = []
for attn_processor_name, attn_processor in unet.attn_processors.items():
# Parse the attention module.
attn_module = unet
for n in attn_processor_name.split(".")[:-1]:
attn_module = getattr(attn_module, n)
# Set the `lora_layer` attribute of the attention-related matrices.
attn_module.to_q.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_q.in_features, out_features=attn_module.to_q.out_features, rank=args.rank
)
)
attn_module.to_k.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_k.in_features, out_features=attn_module.to_k.out_features, rank=args.rank
)
)
attn_module.to_v.set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_v.in_features, out_features=attn_module.to_v.out_features, rank=args.rank
)
)
attn_module.to_out[0].set_lora_layer(
LoRALinearLayer(
in_features=attn_module.to_out[0].in_features,
out_features=attn_module.to_out[0].out_features,
rank=args.rank,
)
)
# Accumulate the LoRA params to optimize.
unet_lora_parameters.extend(attn_module.to_q.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_k.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_v.lora_layer.parameters())
unet_lora_parameters.extend(attn_module.to_out[0].lora_layer.parameters())
unet_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["to_k", "to_q", "to_v", "to_out.0"]
)
unet.add_adapter(unet_lora_config)
# The text encoder comes from 🤗 transformers, so we cannot directly modify it.
# So, instead, we monkey-patch the forward calls of its attention-blocks.
if args.train_text_encoder:
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
text_lora_parameters_one = LoraLoaderMixin._modify_text_encoder(
text_encoder_one, dtype=torch.float32, rank=args.rank
)
text_lora_parameters_two = LoraLoaderMixin._modify_text_encoder(
text_encoder_two, dtype=torch.float32, rank=args.rank
text_lora_config = LoraConfig(
r=args.rank, init_lora_weights="gaussian", target_modules=["q_proj", "k_proj", "v_proj", "out_proj"]
)
text_encoder_one.add_adapter(text_lora_config)
text_encoder_two.add_adapter(text_lora_config)
# create custom saving & loading hooks so that `accelerator.save_state(...)` serializes in a nice format
def save_model_hook(models, weights, output_dir):
@@ -805,11 +1035,11 @@ def main(args):
for model in models:
if isinstance(model, type(accelerator.unwrap_model(unet))):
unet_lora_layers_to_save = unet_lora_state_dict(model)
unet_lora_layers_to_save = get_peft_model_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_one))):
text_encoder_one_lora_layers_to_save = text_encoder_lora_state_dict(model)
text_encoder_one_lora_layers_to_save = get_peft_model_state_dict(model)
elif isinstance(model, type(accelerator.unwrap_model(text_encoder_two))):
text_encoder_two_lora_layers_to_save = text_encoder_lora_state_dict(model)
text_encoder_two_lora_layers_to_save = get_peft_model_state_dict(model)
else:
raise ValueError(f"unexpected save model: {model.__class__}")
@@ -866,35 +1096,109 @@ def main(args):
args.learning_rate * args.gradient_accumulation_steps * args.train_batch_size * accelerator.num_processes
)
# Use 8-bit Adam for lower memory usage or to fine-tune the model in 16GB GPUs
if args.use_8bit_adam:
try:
import bitsandbytes as bnb
except ImportError:
raise ImportError(
"To use 8-bit Adam, please install the bitsandbytes library: `pip install bitsandbytes`."
)
unet_lora_parameters = list(filter(lambda p: p.requires_grad, unet.parameters()))
optimizer_class = bnb.optim.AdamW8bit
if args.train_text_encoder:
text_lora_parameters_one = list(filter(lambda p: p.requires_grad, text_encoder_one.parameters()))
text_lora_parameters_two = list(filter(lambda p: p.requires_grad, text_encoder_two.parameters()))
# Optimization parameters
unet_lora_parameters_with_lr = {"params": unet_lora_parameters, "lr": args.learning_rate}
if args.train_text_encoder:
# different learning rate for text encoder and unet
text_lora_parameters_one_with_lr = {
"params": text_lora_parameters_one,
"weight_decay": args.adam_weight_decay_text_encoder,
"lr": args.text_encoder_lr if args.text_encoder_lr else args.learning_rate,
}
text_lora_parameters_two_with_lr = {
"params": text_lora_parameters_two,
"weight_decay": args.adam_weight_decay_text_encoder,
"lr": args.text_encoder_lr if args.text_encoder_lr else args.learning_rate,
}
params_to_optimize = [
unet_lora_parameters_with_lr,
text_lora_parameters_one_with_lr,
text_lora_parameters_two_with_lr,
]
else:
optimizer_class = torch.optim.AdamW
params_to_optimize = [unet_lora_parameters_with_lr]
# Optimizer creation
params_to_optimize = (
itertools.chain(unet_lora_parameters, text_lora_parameters_one, text_lora_parameters_two)
if args.train_text_encoder
else unet_lora_parameters
if not (args.optimizer.lower() == "prodigy" or args.optimizer.lower() == "adamw"):
logger.warn(
f"Unsupported choice of optimizer: {args.optimizer}.Supported optimizers include [adamW, prodigy]."
"Defaulting to adamW"
)
args.optimizer = "adamw"
if args.use_8bit_adam and not args.optimizer.lower() == "adamw":
logger.warn(
f"use_8bit_adam is ignored when optimizer is not set to 'AdamW'. Optimizer was "
f"set to {args.optimizer.lower()}"
)
if args.optimizer.lower() == "adamw":
if args.use_8bit_adam:
try:
import bitsandbytes as bnb
except ImportError:
raise ImportError(
"To use 8-bit Adam, please install the bitsandbytes library: `pip install bitsandbytes`."
)
optimizer_class = bnb.optim.AdamW8bit
else:
optimizer_class = torch.optim.AdamW
optimizer = optimizer_class(
params_to_optimize,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
)
if args.optimizer.lower() == "prodigy":
try:
import prodigyopt
except ImportError:
raise ImportError("To use Prodigy, please install the prodigyopt library: `pip install prodigyopt`")
optimizer_class = prodigyopt.Prodigy
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
decouple=args.prodigy_decouple,
use_bias_correction=args.prodigy_use_bias_correction,
safeguard_warmup=args.prodigy_safeguard_warmup,
)
# Dataset and DataLoaders creation:
train_dataset = DreamBoothDataset(
instance_data_root=args.instance_data_dir,
instance_prompt=args.instance_prompt,
class_prompt=args.class_prompt,
class_data_root=args.class_data_dir if args.with_prior_preservation else None,
class_num=args.num_class_images,
size=args.resolution,
repeats=args.repeats,
center_crop=args.center_crop,
)
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
train_dataloader = torch.utils.data.DataLoader(
train_dataset,
batch_size=args.train_batch_size,
shuffle=True,
collate_fn=lambda examples: collate_fn(examples, args.with_prior_preservation),
num_workers=args.dataloader_num_workers,
)
# Computes additional embeddings/ids required by the SDXL UNet.
# regular text emebddings (when `train_text_encoder` is not True)
# regular text embeddings (when `train_text_encoder` is not True)
# pooled text embeddings
# time ids
@@ -921,7 +1225,11 @@ def main(args):
# Handle instance prompt.
instance_time_ids = compute_time_ids()
if not args.train_text_encoder:
# If no type of tuning is done on the text_encoder and custom instance prompts are NOT
# provided (i.e. the --instance_prompt is used for all images), we encode the instance prompt once to avoid
# the redundant encoding.
if not args.train_text_encoder and not train_dataset.custom_instance_prompts:
instance_prompt_hidden_states, instance_pooled_prompt_embeds = compute_text_embeddings(
args.instance_prompt, text_encoders, tokenizers
)
@@ -934,49 +1242,36 @@ def main(args):
args.class_prompt, text_encoders, tokenizers
)
# Clear the memory here.
if not args.train_text_encoder:
# Clear the memory here
if not args.train_text_encoder and not train_dataset.custom_instance_prompts:
del tokenizers, text_encoders
gc.collect()
torch.cuda.empty_cache()
# Pack the statically computed variables appropriately. This is so that we don't
# If custom instance prompts are NOT provided (i.e. the instance prompt is used for all images),
# pack the statically computed variables appropriately here. This is so that we don't
# have to pass them to the dataloader.
add_time_ids = instance_time_ids
if args.with_prior_preservation:
add_time_ids = torch.cat([add_time_ids, class_time_ids], dim=0)
if not args.train_text_encoder:
prompt_embeds = instance_prompt_hidden_states
unet_add_text_embeds = instance_pooled_prompt_embeds
if args.with_prior_preservation:
prompt_embeds = torch.cat([prompt_embeds, class_prompt_hidden_states], dim=0)
unet_add_text_embeds = torch.cat([unet_add_text_embeds, class_pooled_prompt_embeds], dim=0)
else:
tokens_one = tokenize_prompt(tokenizer_one, args.instance_prompt)
tokens_two = tokenize_prompt(tokenizer_two, args.instance_prompt)
if args.with_prior_preservation:
class_tokens_one = tokenize_prompt(tokenizer_one, args.class_prompt)
class_tokens_two = tokenize_prompt(tokenizer_two, args.class_prompt)
tokens_one = torch.cat([tokens_one, class_tokens_one], dim=0)
tokens_two = torch.cat([tokens_two, class_tokens_two], dim=0)
# Dataset and DataLoaders creation:
train_dataset = DreamBoothDataset(
instance_data_root=args.instance_data_dir,
class_data_root=args.class_data_dir if args.with_prior_preservation else None,
class_num=args.num_class_images,
size=args.resolution,
center_crop=args.center_crop,
)
train_dataloader = torch.utils.data.DataLoader(
train_dataset,
batch_size=args.train_batch_size,
shuffle=True,
collate_fn=lambda examples: collate_fn(examples, args.with_prior_preservation),
num_workers=args.dataloader_num_workers,
)
if not train_dataset.custom_instance_prompts:
if not args.train_text_encoder:
prompt_embeds = instance_prompt_hidden_states
unet_add_text_embeds = instance_pooled_prompt_embeds
if args.with_prior_preservation:
prompt_embeds = torch.cat([prompt_embeds, class_prompt_hidden_states], dim=0)
unet_add_text_embeds = torch.cat([unet_add_text_embeds, class_pooled_prompt_embeds], dim=0)
# if we're optmizing the text encoder (both if instance prompt is used for all images or custom prompts) we need to tokenize and encode the
# batch prompts on all training steps
else:
tokens_one = tokenize_prompt(tokenizer_one, args.instance_prompt)
tokens_two = tokenize_prompt(tokenizer_two, args.instance_prompt)
if args.with_prior_preservation:
class_tokens_one = tokenize_prompt(tokenizer_one, args.class_prompt)
class_tokens_two = tokenize_prompt(tokenizer_two, args.class_prompt)
tokens_one = torch.cat([tokens_one, class_tokens_one], dim=0)
tokens_two = torch.cat([tokens_two, class_tokens_two], dim=0)
# Scheduler and math around the number of training steps.
overrode_max_train_steps = False
@@ -1079,6 +1374,17 @@ def main(args):
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(unet):
pixel_values = batch["pixel_values"].to(dtype=vae.dtype)
prompts = batch["prompts"]
# encode batch prompts when custom prompts are provided for each image -
if train_dataset.custom_instance_prompts:
if not args.train_text_encoder:
prompt_embeds, unet_add_text_embeds = compute_text_embeddings(
prompts, text_encoders, tokenizers
)
else:
tokens_one = tokenize_prompt(tokenizer_one, prompts)
tokens_two = tokenize_prompt(tokenizer_two, prompts)
# Convert images to latent space
model_input = vae.encode(pixel_values).latent_dist.sample()
@@ -1099,16 +1405,21 @@ def main(args):
# (this is the forward diffusion process)
noisy_model_input = noise_scheduler.add_noise(model_input, noise, timesteps)
# Calculate the elements to repeat depending on the use of prior-preservation.
elems_to_repeat = bsz // 2 if args.with_prior_preservation else bsz
# Calculate the elements to repeat depending on the use of prior-preservation and custom captions.
if not train_dataset.custom_instance_prompts:
elems_to_repeat_text_embeds = bsz // 2 if args.with_prior_preservation else bsz
elems_to_repeat_time_ids = bsz // 2 if args.with_prior_preservation else bsz
else:
elems_to_repeat_text_embeds = 1
elems_to_repeat_time_ids = bsz // 2 if args.with_prior_preservation else bsz
# Predict the noise residual
if not args.train_text_encoder:
unet_added_conditions = {
"time_ids": add_time_ids.repeat(elems_to_repeat, 1),
"text_embeds": unet_add_text_embeds.repeat(elems_to_repeat, 1),
"time_ids": add_time_ids.repeat(elems_to_repeat_time_ids, 1),
"text_embeds": unet_add_text_embeds.repeat(elems_to_repeat_text_embeds, 1),
}
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat, 1, 1)
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat_text_embeds, 1, 1)
model_pred = unet(
noisy_model_input,
timesteps,
@@ -1116,15 +1427,17 @@ def main(args):
added_cond_kwargs=unet_added_conditions,
).sample
else:
unet_added_conditions = {"time_ids": add_time_ids.repeat(elems_to_repeat, 1)}
unet_added_conditions = {"time_ids": add_time_ids.repeat(elems_to_repeat_time_ids, 1)}
prompt_embeds, pooled_prompt_embeds = encode_prompt(
text_encoders=[text_encoder_one, text_encoder_two],
tokenizers=None,
prompt=None,
text_input_ids_list=[tokens_one, tokens_two],
)
unet_added_conditions.update({"text_embeds": pooled_prompt_embeds.repeat(elems_to_repeat, 1)})
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat, 1, 1)
unet_added_conditions.update(
{"text_embeds": pooled_prompt_embeds.repeat(elems_to_repeat_text_embeds, 1)}
)
prompt_embeds_input = prompt_embeds.repeat(elems_to_repeat_text_embeds, 1, 1)
model_pred = unet(
noisy_model_input, timesteps, prompt_embeds_input, added_cond_kwargs=unet_added_conditions
).sample
@@ -1142,16 +1455,34 @@ def main(args):
model_pred, model_pred_prior = torch.chunk(model_pred, 2, dim=0)
target, target_prior = torch.chunk(target, 2, dim=0)
# Compute instance loss
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
# Compute prior loss
prior_loss = F.mse_loss(model_pred_prior.float(), target_prior.float(), reduction="mean")
if args.snr_gamma is None:
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
else:
# Compute loss-weights as per Section 3.4 of https://arxiv.org/abs/2303.09556.
# Since we predict the noise instead of x_0, the original formulation is slightly changed.
# This is discussed in Section 4.2 of the same paper.
snr = compute_snr(noise_scheduler, timesteps)
base_weight = (
torch.stack([snr, args.snr_gamma * torch.ones_like(timesteps)], dim=1).min(dim=1)[0] / snr
)
if noise_scheduler.config.prediction_type == "v_prediction":
# Velocity objective needs to be floored to an SNR weight of one.
mse_loss_weights = base_weight + 1
else:
# Epsilon and sample both use the same loss weights.
mse_loss_weights = base_weight
loss = F.mse_loss(model_pred.float(), target.float(), reduction="none")
loss = loss.mean(dim=list(range(1, len(loss.shape)))) * mse_loss_weights
loss = loss.mean()
if args.with_prior_preservation:
# Add the prior loss to the instance loss.
loss = loss + args.prior_loss_weight * prior_loss
else:
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
accelerator.backward(loss)
if accelerator.sync_gradients:
@@ -1212,10 +1543,16 @@ def main(args):
# create pipeline
if not args.train_text_encoder:
text_encoder_one = text_encoder_cls_one.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path,
subfolder="text_encoder",
revision=args.revision,
variant=args.variant,
)
text_encoder_two = text_encoder_cls_two.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision
args.pretrained_model_name_or_path,
subfolder="text_encoder_2",
revision=args.revision,
variant=args.variant,
)
pipeline = StableDiffusionXLPipeline.from_pretrained(
args.pretrained_model_name_or_path,
@@ -1224,6 +1561,7 @@ def main(args):
text_encoder_2=accelerator.unwrap_model(text_encoder_two),
unet=accelerator.unwrap_model(unet),
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
@@ -1277,13 +1615,13 @@ def main(args):
if accelerator.is_main_process:
unet = accelerator.unwrap_model(unet)
unet = unet.to(torch.float32)
unet_lora_layers = unet_lora_state_dict(unet)
unet_lora_layers = get_peft_model_state_dict(unet)
if args.train_text_encoder:
text_encoder_one = accelerator.unwrap_model(text_encoder_one)
text_encoder_lora_layers = text_encoder_lora_state_dict(text_encoder_one.to(torch.float32))
text_encoder_lora_layers = get_peft_model_state_dict(text_encoder_one.to(torch.float32))
text_encoder_two = accelerator.unwrap_model(text_encoder_two)
text_encoder_2_lora_layers = text_encoder_lora_state_dict(text_encoder_two.to(torch.float32))
text_encoder_2_lora_layers = get_peft_model_state_dict(text_encoder_two.to(torch.float32))
else:
text_encoder_lora_layers = None
text_encoder_2_lora_layers = None
@@ -1301,10 +1639,15 @@ def main(args):
vae_path,
subfolder="vae" if args.pretrained_vae_model_name_or_path is None else None,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline = StableDiffusionXLPipeline.from_pretrained(
args.pretrained_model_name_or_path, vae=vae, revision=args.revision, torch_dtype=weight_dtype
args.pretrained_model_name_or_path,
vae=vae,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
# We train on the simplified learning objective. If we were previously predicting a variance, we need the scheduler to ignore it
@@ -1353,7 +1696,8 @@ def main(args):
images=images,
base_model=args.pretrained_model_name_or_path,
train_text_encoder=args.train_text_encoder,
prompt=args.instance_prompt,
instance_prompt=args.instance_prompt,
validation_prompt=args.validation_prompt,
repo_folder=args.output_dir,
vae_path=args.pretrained_vae_model_name_or_path,
)
@@ -0,0 +1,101 @@
# coding=utf-8
# Copyright 2023 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import sys
import tempfile
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class InstructPix2Pix(ExamplesTestsAccelerate):
def test_instruct_pix2pix_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/instruct_pix2pix/train_instruct_pix2pix.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/instructpix2pix-10-samples
--resolution=64
--random_flip
--train_batch_size=1
--max_train_steps=7
--checkpointing_steps=2
--checkpoints_total_limit=2
--output_dir {tmpdir}
--seed=0
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_instruct_pix2pix_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
examples/instruct_pix2pix/train_instruct_pix2pix.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/instructpix2pix-10-samples
--resolution=64
--random_flip
--train_batch_size=1
--max_train_steps=9
--checkpointing_steps=2
--output_dir {tmpdir}
--seed=0
""".split()
run_command(self._launch_args + test_args)
# check checkpoint directories exist
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4", "checkpoint-6", "checkpoint-8"},
)
resume_run_args = f"""
examples/instruct_pix2pix/train_instruct_pix2pix.py
--pretrained_model_name_or_path=hf-internal-testing/tiny-stable-diffusion-pipe
--dataset_name=hf-internal-testing/instructpix2pix-10-samples
--resolution=64
--random_flip
--train_batch_size=1
--max_train_steps=11
--checkpointing_steps=2
--output_dir {tmpdir}
--seed=0
--resume_from_checkpoint=checkpoint-8
--checkpoints_total_limit=3
""".split()
run_command(self._launch_args + resume_run_args)
# check checkpoint directories exist
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-6", "checkpoint-8", "checkpoint-10"},
)
@@ -52,7 +52,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -78,6 +78,12 @@ def parse_args():
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--dataset_name",
type=str,
@@ -435,9 +441,11 @@ def main():
args.pretrained_model_name_or_path, subfolder="tokenizer", revision=args.revision
)
text_encoder = CLIPTextModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision, variant=args.variant
)
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae", revision=args.revision)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.non_ema_revision
)
@@ -915,6 +923,7 @@ def main():
text_encoder=accelerator.unwrap_model(text_encoder),
vae=accelerator.unwrap_model(vae),
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline = pipeline.to(accelerator.device)
@@ -966,6 +975,7 @@ def main():
vae=accelerator.unwrap_model(vae),
unet=unet,
revision=args.revision,
variant=args.variant,
)
pipeline.save_pretrained(args.output_dir)
@@ -55,7 +55,7 @@ from diffusers.utils.import_utils import is_xformers_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -118,6 +118,12 @@ def parse_args():
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--variant",
type=str,
default=None,
help="Variant of the model files of the pretrained model identifier from huggingface.co/models, 'e.g.' fp16",
)
parser.add_argument(
"--dataset_name",
type=str,
@@ -484,9 +490,10 @@ def main():
vae_path,
subfolder="vae" if args.pretrained_vae_model_name_or_path is None else None,
revision=args.revision,
variant=args.variant,
)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, variant=args.variant
)
# InstructPix2Pix uses an additional image for conditioning. To accommodate that,
@@ -695,10 +702,16 @@ def main():
# Load scheduler, tokenizer and models.
tokenizer_1 = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer",
revision=args.revision,
use_fast=False,
)
tokenizer_2 = AutoTokenizer.from_pretrained(
args.pretrained_model_name_or_path, subfolder="tokenizer_2", revision=args.revision, use_fast=False
args.pretrained_model_name_or_path,
subfolder="tokenizer_2",
revision=args.revision,
use_fast=False,
)
text_encoder_cls_1 = import_model_class_from_model_name_or_path(args.pretrained_model_name_or_path, args.revision)
text_encoder_cls_2 = import_model_class_from_model_name_or_path(
@@ -708,10 +721,10 @@ def main():
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder_1 = text_encoder_cls_1.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder", revision=args.revision, variant=args.variant
)
text_encoder_2 = text_encoder_cls_2.from_pretrained(
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision
args.pretrained_model_name_or_path, subfolder="text_encoder_2", revision=args.revision, variant=args.variant
)
# We ALWAYS pre-compute the additional condition embeddings needed for SDXL
@@ -1109,6 +1122,7 @@ def main():
tokenizer_2=tokenizer_2,
vae=vae,
revision=args.revision,
variant=args.variant,
torch_dtype=weight_dtype,
)
pipeline = pipeline.to(accelerator.device)
@@ -1176,6 +1190,7 @@ def main():
vae=vae,
unet=unet,
revision=args.revision,
variant=args.variant,
)
pipeline.save_pretrained(args.output_dir)
@@ -52,7 +52,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -46,7 +46,7 @@ from diffusers.utils import check_min_version, is_wandb_available
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
@@ -51,7 +51,7 @@ if is_wandb_available():
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.24.0.dev0")
check_min_version("0.25.0.dev0")
logger = get_logger(__name__, log_level="INFO")
+2 -2
View File
@@ -1,7 +1,7 @@
# Research projects
This folder contains various research projects using 🧨 Diffusers.
They are not really maintained by the core maintainers of this library and often require a specific version of Diffusers that is indicated in the requirements file of each folder.
This folder contains various research projects using 🧨 Diffusers.
They are not really maintained by the core maintainers of this library and often require a specific version of Diffusers that is indicated in the requirements file of each folder.
Updating them to the most recent version of the library will require some work.
To use any of them, just run the command
@@ -420,7 +420,7 @@ def import_model_class_from_model_name_or_path(
pretrained_model_name_or_path: str, revision: str, subfolder: str = "text_encoder"
):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path, subfolder=subfolder, revision=revision, use_auth_token=True
pretrained_model_name_or_path, subfolder=subfolder, revision=revision
)
model_class = text_encoder_config.architectures[0]
@@ -975,7 +975,7 @@ def main(args):
revision=args.revision,
)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision, use_auth_token=True
args.pretrained_model_name_or_path, subfolder="unet", revision=args.revision
)
if args.controlnet_model_name_or_path:
@@ -1,6 +1,6 @@
## [Deprecated] Multi Token Textual Inversion
**IMPORTART: This research project is deprecated. Multi Token Textual Inversion is now supported natively in [the officail textual inversion example](https://github.com/huggingface/diffusers/tree/main/examples/textual_inversion#running-locally-with-pytorch).**
**IMPORTART: This research project is deprecated. Multi Token Textual Inversion is now supported natively in [the official textual inversion example](https://github.com/huggingface/diffusers/tree/main/examples/textual_inversion#running-locally-with-pytorch).**
The author of this project is [Isamu Isozaki](https://github.com/isamu-isozaki) - please make sure to tag the author for issue and PRs as well as @patrickvonplaten.
@@ -17,9 +17,9 @@ Feel free to add these options to your training! In practice num_vec_per_token a
[Textual inversion](https://arxiv.org/abs/2208.01618) is a method to personalize text2image models like stable diffusion on your own images using just 3-5 examples.
The `textual_inversion.py` script shows how to implement the training procedure and adapt it for stable diffusion.
## Running on Colab
## Running on Colab
Colab for training
Colab for training
[![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_textual_inversion_training.ipynb)
Colab for inference
@@ -53,7 +53,7 @@ accelerate config
### Cat toy example
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-5`, so you'll need to visit [its card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license and tick the checkbox if you agree.
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-5`, so you'll need to visit [its card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license and tick the checkbox if you agree.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
@@ -63,7 +63,7 @@ Run the following command to authenticate your token
huggingface-cli login
```
If you have already cloned the repo, then you won't need to go through these steps.
If you have already cloned the repo, then you won't need to go through these steps.
<br>
@@ -2,4 +2,4 @@
**This research project is not actively maintained by the diffusers team. For any questions or comments, please contact Prathik Rao (prathikr), Sunghoon Choi (hanbitmyths), Ashwini Khade (askhade), or Peng Wang (pengwa) on github with any questions.**
This aims to provide diffusers examples with ONNXRuntime optimizations for training/fine-tuning unconditional image generation, text to image, and textual inversion. Please see individual directories for more details on how to run each task using ONNXRuntime.
This aims to provide diffusers examples with ONNXRuntime optimizations for training/fine-tuning unconditional image generation, text to image, and textual inversion. Please see individual directories for more details on how to run each task using ONNXRuntime.
@@ -34,7 +34,7 @@ accelerate config
### Pokemon example
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-4`, so you'll need to visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4), read the license and tick the checkbox if you agree.
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-4`, so you'll need to visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4), read the license and tick the checkbox if you agree.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
@@ -68,7 +68,7 @@ accelerate launch --mixed_precision="fp16" train_text_to_image.py \
--learning_rate=1e-05 \
--max_grad_norm=1 \
--lr_scheduler="constant" --lr_warmup_steps=0 \
--output_dir="sd-pokemon-model"
--output_dir="sd-pokemon-model"
```
Please contact Prathik Rao (prathikr), Sunghoon Choi (hanbitmyths), Ashwini Khade (askhade), or Peng Wang (pengwa) on github with any questions.
@@ -3,9 +3,9 @@
[Textual inversion](https://arxiv.org/abs/2208.01618) is a method to personalize text2image models like stable diffusion on your own images using just 3-5 examples.
The `textual_inversion.py` script shows how to implement the training procedure and adapt it for stable diffusion.
## Running on Colab
## Running on Colab
Colab for training
Colab for training
[![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_textual_inversion_training.ipynb)
Colab for inference
@@ -39,7 +39,7 @@ accelerate config
### Cat toy example
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-5`, so you'll need to visit [its card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license and tick the checkbox if you agree.
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-5`, so you'll need to visit [its card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license and tick the checkbox if you agree.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
@@ -49,7 +49,7 @@ Run the following command to authenticate your token
huggingface-cli login
```
If you have already cloned the repo, then you won't need to go through these steps.
If you have already cloned the repo, then you won't need to go through these steps.
<br>

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