* modulize log validation
* run make style and refactor wanddb support
* remove redundant initialization
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* make checkpoint_merger pipeline pass the "variant" argument to from_pretrained()
* make style
---------
Co-authored-by: Lincoln Stein <lstein@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add stable_diffusion_xl_ipex community pipeline
* make style for code quality check
* update docs as suggested
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* standardize model card
* fix tags
* correct import styling and update tags
* run make style and make quality
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: allow low_cpu_mem_usage in ip adapter loading
* reduce the number of device placements.
* documentation.
* throw low_cpu_mem_usage warning only once from the main entry point.
* use load_model_into_meta in single file utils
* propagate to autoencoder and controlnet.
* correct class name access behaviour.
* remove torch_dtype from load_model_into_meta; seems unncessary
* remove incorrect kwarg
* style to avoid extra unnecessary line breaks
* fix: bias loading bug
* fixes for SDXL
* apply changes to the conversion script to match single_file_utils.py
* do transpose to match the single file loading logic.
Remove <cat-toy> validation prompt from textual_inversion_sdxl.py
The `<cat-toy>` validation prompt is a default choice for the example task in the README. But no other part of `textual_inversion_sdxl.py` references the cat toy and `textual_inversion.py` has a default validation prompt of `None` as well.
So bring `textual_inversion_sdxl.py` in line with `textual_inversion.py` and change default validation prompt to `None`
* attention_head_dim
* debug
* print more info
* correct num_attention_heads behaviour
* down_block_num_attention_heads -> num_attention_heads.
* correct the image link in doc.
* add: deprecation for num_attention_head
* fix: test argument to use attention_head_dim
* more fixes.
* quality
* address comments.
* remove depcrecation.
* add: support for passing ip adapter image embeddings
* debugging
* make feature_extractor unloading conditioned on safety_checker
* better condition
* type annotation
* index to look into value slices
* more debugging
* debugging
* serialize embeddings dict
* better conditioning
* remove unnecessary prints.
* Update src/diffusers/loaders/ip_adapter.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* make fix-copies and styling.
* styling and further copy fixing.
* fix: check_inputs call in controlnet sdxl img2img pipeline
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* feat: standarize model card creation for dreambooth training.
* correct 'inference
* remove comments.
* take component out of kwargs
* style
* add: card template to have a leaner description.
* widget support.
* propagate changes to train_dreambooth_lora
* propagate changes to custom diffusion
* make widget properly type-annotated
* fix: callback function name is incorrect
On this tutorial there is a function defined and then used inside `callback_on_step_end` argument, but the name was not correct (mismatch)
* fix: typo in num_timestep (correct is num_timesteps)
fixed property name
* remove _to_tensor
* remove _to_tensor definition
* remove _collapse_frames_into_batch
* remove lora for not bloating the code.
* remove sample_size.
* simplify code a bit more
* ensure timesteps are always in tensor.
* Fix `AutoencoderTiny` with `use_slicing`
When using slicing with AutoencoderTiny, the encoder mistakenly encodes the entire batch for every image in the batch.
* Fixed formatting issue
* add noise_offset param
* micro conditioning - wip
* image processing adjusted and moved to support micro conditioning
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* time ids shape fix
* move token replacement of validation prompt to the same section of instance prompt and class prompt
* add offset noise to sd15 advanced script
* fix token loading during validation
* fix token loading during validation in sdxl script
* a little clean
* style
* a little clean
* style
* sdxl script - a little clean + minor path fix
sd 1.5 script - change default resolution value
* ad 1.5 script - minor path fix
* fix missing comma in code example in model card
* clean up commented lines
* style
* remove time ids computed outside training loop - no longer used now that we utilize micro-conditioning, as all time ids are now computed inside the training loop
* style
* [WIP] - added draft readme, building off of examples/dreambooth/README.md
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* removed --crops_coords_top_left from CLI args
* style
* fix missing shape bug due to missing RGB if statement
* add blog mention at the start of the reamde as well
* Update examples/advanced_diffusion_training/README.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* change note to render nicely as well
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Fix bug in ResnetBlock2D.forward when not USE_PEFT_BACKEND and using scale_shift for time emb where the lora scale gets overwritten.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix minsnr implementation for v-prediction case
* format code
* always compute snr when snr_gamma is specified
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: explicitly tag to diffusers when using push_to_hub
* remove tags.
* reset repo.
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix: tests
* fix: push_to_hub behaviour for tagging from save_pretrained
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* import fixes.
* add library name to existing model card.
* add: standalone test for generate_model_card
* fix tests for standalone method
* moved library_name to a better place.
* merge create_model_card and generate_model_card.
* fix test
* address lucain's comments
* fix return identation
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* address further comments.
* Update src/diffusers/pipelines/pipeline_utils.py
Co-authored-by: Lucain <lucainp@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Lucain <lucainp@gmail.com>
* initial commit for unconditional/class-conditional consistency training script
* make style
* Add entry for consistency training script in community README.
* Move consistency training script from community to research_projects/consistency_training
* Add requirements.txt and README to research_projects/consistency_training directory.
* Manually revert community README changes for consistency training.
* Fix path to script after moving script to research projects.
* Add option to load U-Net weights from pretrained model.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* begin animatediff img2video and video2video
* revert animatediff to original implementation
* add img2video as pipeline
* update
* add vid2vid pipeline
* update imports
* update
* remove copied from line for check_inputs
* update
* update examples
* add multi-batch support
* fix __init__.py files
* move img2vid to community
* update community readme and examples
* fix
* make fix-copies
* add vid2vid batch params
* apply suggestions from review
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* add test for animatediff vid2vid
* torch.stack -> torch.cat
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* make style
* docs for vid2vid
* update
* fix prepare_latents
* fix docs
* remove img2vid
* update README to :main
* remove slow test
* refactor pipeline output
* update docs
* update docs
* merge community readme from :main
* final fix i promise
* add support for url in animatediff example
* update example
* update callbacks to latest implementation
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix merge
* Apply suggestions from code review
* remove callback and callback_steps as suggested in review
* Update tests/pipelines/animatediff/test_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix import error caused due to unet refactor in #6630
* fix numpy import error after tensor2vid refactor in #6626
* make fix-copies
* fix numpy error
* fix progress bar test
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* sd1.5 support in separate script
A quick adaptation to support people interested in using this method on 1.5 models.
* sd15 prompt text encoding and unet conversions
as per @linoytsaban 's recommendations. Testing would be appreciated,
* Readability and quality improvements
Removed some mentions of SDXL, and some arguments that don't apply to sd 1.5, and cleaned up some comments.
* make style/quality commands
* tracker rename and run-it doc
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
---------
Co-authored-by: Linoy Tsaban <57615435+linoytsaban@users.noreply.github.com>
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,11 +12,11 @@ specific language governing permissions and limitations under the License.
# IP-Adapter
[IP-Adapter](https://hf.co/papers/2308.06721) is a lightweight adapter that enables prompting a diffusion model with an image. This method decouples the cross-attention layers of the image and text features. The image features are generated from an image encoder. Files generated from IP-Adapter are only ~100MBs.
[IP-Adapter](https://hf.co/papers/2308.06721) is a lightweight adapter that enables prompting a diffusion model with an image. This method decouples the cross-attention layers of the image and text features. The image features are generated from an image encoder.
<Tip>
Learn how to load an IP-Adapter checkpoint and image in the [IP-Adapter](../../using-diffusers/loading_adapters#ip-adapter) loading guide.
Learn how to load an IP-Adapter checkpoint and image in the IP-Adapter [loading](../../using-diffusers/loading_adapters#ip-adapter) guide, and you can see how to use it in the [usage](../../using-diffusers/ip_adapter) guide.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# UVit2DModel
The [U-ViT](https://hf.co/papers/2301.11093) model is a vision transformer (ViT) based UNet. This model incorporates elements from ViT (considers all inputs such as time, conditions and noisy image patches as tokens) and a UNet (long skip connections between the shallow and deep layers). The skip connection is important for predicting pixel-level features. An additional 3x3 convolutional block is applied prior to the final output to improve image quality.
The abstract from the paper is:
*Currently, applying diffusion models in pixel space of high resolution images is difficult. Instead, existing approaches focus on diffusion in lower dimensional spaces (latent diffusion), or have multiple super-resolution levels of generation referred to as cascades. The downside is that these approaches add additional complexity to the diffusion framework. This paper aims to improve denoising diffusion for high resolution images while keeping the model as simple as possible. The paper is centered around the research question: How can one train a standard denoising diffusion models on high resolution images, and still obtain performance comparable to these alternate approaches? The four main findings are: 1) the noise schedule should be adjusted for high resolution images, 2) It is sufficient to scale only a particular part of the architecture, 3) dropout should be added at specific locations in the architecture, and 4) downsampling is an effective strategy to avoid high resolution feature maps. Combining these simple yet effective techniques, we achieve state-of-the-art on image generation among diffusion models without sampling modifiers on ImageNet.*
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -25,6 +25,7 @@ The abstract of the paper is the following:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py) | *Video-to-Video Generation with AnimateDiff* |
## Available checkpoints
@@ -32,6 +33,8 @@ Motion Adapter checkpoints can be found under [guoyww](https://huggingface.co/gu
## Usage example
### AnimateDiffPipeline
AnimateDiff works with a MotionAdapter checkpoint and a Stable Diffusion model checkpoint. The MotionAdapter is a collection of Motion Modules that are responsible for adding coherent motion across image frames. These modules are applied after the Resnet and Attention blocks in Stable Diffusion UNet.
The following example demonstrates how to use a *MotionAdapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
@@ -98,6 +101,114 @@ AnimateDiff tends to work better with finetuned Stable Diffusion models. If you
</Tip>
### AnimateDiffVideoToVideoPipeline
AnimateDiff can also be used to generate visually similar videos or enable style/character/background or other edits starting from an initial video, allowing you to seamlessly explore creative possibilities.
alt="closeup of tony stark, robert downey jr, fireworks"
style="width: 300px;" />
</td>
</tr>
</table>
## Using Motion LoRAs
Motion LoRAs are a collection of LoRAs that work with the `guoyww/animatediff-motion-adapter-v1-5-2` checkpoint. These LoRAs are responsible for adding specific types of motion to the animations.
@@ -297,19 +408,102 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
</Tip>
## Using AnimateLCM
[AnimateLCM](https://animatelcm.github.io/) is a motion module checkpoint and an [LCM LoRA](https://huggingface.co/docs/diffusers/using-diffusers/inference_with_lcm_lora) that have been created using a consistency learning strategy that decouples the distillation of the image generation priors and the motion generation priors.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# I2VGen-XL
[I2VGen-XL: High-Quality Image-to-Video Synthesis via Cascaded Diffusion Models](https://hf.co/papers/2311.04145.pdf) by Shiwei Zhang, Jiayu Wang, Yingya Zhang, Kang Zhao, Hangjie Yuan, Zhiwu Qin, Xiang Wang, Deli Zhao, and Jingren Zhou.
The abstract from the paper is:
*Video synthesis has recently made remarkable strides benefiting from the rapid development of diffusion models. However, it still encounters challenges in terms of semantic accuracy, clarity and spatio-temporal continuity. They primarily arise from the scarcity of well-aligned text-video data and the complex inherent structure of videos, making it difficult for the model to simultaneously ensure semantic and qualitative excellence. In this report, we propose a cascaded I2VGen-XL approach that enhances model performance by decoupling these two factors and ensures the alignment of the input data by utilizing static images as a form of crucial guidance. I2VGen-XL consists of two stages: i) the base stage guarantees coherent semantics and preserves content from input images by using two hierarchical encoders, and ii) the refinement stage enhances the video's details by incorporating an additional brief text and improves the resolution to 1280×720. To improve the diversity, we collect around 35 million single-shot text-video pairs and 6 billion text-image pairs to optimize the model. By this means, I2VGen-XL can simultaneously enhance the semantic accuracy, continuity of details and clarity of generated videos. Through extensive experiments, we have investigated the underlying principles of I2VGen-XL and compared it with current top methods, which can demonstrate its effectiveness on diverse data. The source code and models will be publicly available at [this https URL](https://i2vgen-xl.github.io/).*
The original codebase can be found [here](https://github.com/ali-vilab/i2vgen-xl/). The model checkpoints can be found [here](https://huggingface.co/ali-vilab/).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Video Diffusion
Stable Video Diffusion was proposed in [Stable Video Diffusion: Scaling Latent Video Diffusion Models to Large Datasets](https://hf.co/papers/2311.15127) by Andreas Blattmann, Tim Dockhorn, Sumith Kulal, Daniel Mendelevitch, Maciej Kilian, Dominik Lorenz, Yam Levi, Zion English, Vikram Voleti, Adam Letts, Varun Jampani, Robin Rombach.
The abstract from the paper is:
*We present Stable Video Diffusion - a latent video diffusion model for high-resolution, state-of-the-art text-to-video and image-to-video generation. Recently, latent diffusion models trained for 2D image synthesis have been turned into generative video models by inserting temporal layers and finetuning them on small, high-quality video datasets. However, training methods in the literature vary widely, and the field has yet to agree on a unified strategy for curating video data. In this paper, we identify and evaluate three different stages for successful training of video LDMs: text-to-image pretraining, video pretraining, and high-quality video finetuning. Furthermore, we demonstrate the necessity of a well-curated pretraining dataset for generating high-quality videos and present a systematic curation process to train a strong base model, including captioning and filtering strategies. We then explore the impact of finetuning our base model on high-quality data and train a text-to-video model that is competitive with closed-source video generation. We also show that our base model provides a powerful motion representation for downstream tasks such as image-to-video generation and adaptability to camera motion-specific LoRA modules. Finally, we demonstrate that our model provides a strong multi-view 3D-prior and can serve as a base to finetune a multi-view diffusion model that jointly generates multiple views of objects in a feedforward fashion, outperforming image-based methods at a fraction of their compute budget. We release code and model weights at this https URL.*
<Tip>
To learn how to use Stable Video Diffusion, take a look at the [Stable Video Diffusion](../../../using-diffusers/svd) guide.
<br>
Check out the [Stability AI](https://huggingface.co/stabilityai) Hub organization for the [base](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid) and [extended frame](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid-xt) checkpoints!
</Tip>
## Tips
Video generation is memory-intensive and one way to reduce your memory usage is to set `enable_forward_chunking` on the pipeline's UNet so you don't run the entire feedforward layer at once. Breaking it up into chunks in a loop is more efficient.
Check out the [Text or image-to-video](text-img2vid) guide for more details about how certain parameters can affect video generation and how to optimize inference by reducing memory usage.
@@ -167,6 +167,12 @@ Here are some sample outputs:
</tr>
</table>
## Tips
Video generation is memory-intensive and one way to reduce your memory usage is to set `enable_forward_chunking` on the pipeline's UNet so you don't run the entire feedforward layer at once. Breaking it up into chunks in a loop is more efficient.
Check out the [Text or image-to-video](text-img2vid) guide for more details about how certain parameters can affect video generation and how to optimize inference by reducing memory usage.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
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