[wip: doc-builder test]

.github/ISSUE_TEMPLATE/bug-report.yml

@@ -49,32 +49,3 @@ body:
       placeholder: diffusers version, platform, python version, ...
     validations:
       required: true
-  - type: textarea
-    id: who-can-help
-    attributes:
-      label: Who can help?
-      description: |
-        Your issue will be replied to more quickly if you can figure out the right person to tag with @
-        If you know how to use git blame, that is the easiest way, otherwise, here is a rough guide of **who to tag**.
-        
-        All issues are read by one of the core maintainers, so if you don't know who to tag, just leave this blank and
-        a core maintainer will ping the right person.
-        
-        Please tag fewer than 3 people.
-        
-        General library related questions: @patrickvonplaten and @sayakpaul
-
-        Questions on the training examples: @williamberman, @sayakpaul, @yiyixuxu
-
-        Questions on memory optimizations, LoRA, float16, etc.: @williamberman, @patrickvonplaten, and @sayakpaul
-
-        Questions on schedulers: @patrickvonplaten and @williamberman
-
-        Questions on models and pipelines: @patrickvonplaten, @sayakpaul, and @williamberman
-
-        Questions on JAX- and MPS-related things: @pcuenca
-
-        Questions on audio pipelines: @patrickvonplaten, @kashif, and @sanchit-gandhi 
-        
-        Documentation: @stevhliu and @yiyixuxu
-      placeholder: "@Username ..."

.github/PULL_REQUEST_TEMPLATE.md

@@ -1,60 +0,0 @@
-# What does this PR do?
-
-<!--
-Congratulations! You've made it this far! You're not quite done yet though.
-
-Once merged, your PR is going to appear in the release notes with the title you set, so make sure it's a great title that fully reflects the extent of your awesome contribution.
-
-Then, please replace this with a description of the change and which issue is fixed (if applicable). Please also include relevant motivation and context. List any dependencies (if any) that are required for this change.
-
-Once you're done, someone will review your PR shortly (see the section "Who can review?" below to tag some potential reviewers). They may suggest changes to make the code even better. If no one reviewed your PR after a week has passed, don't hesitate to post a new comment @-mentioning the same persons---sometimes notifications get lost.
--->
-
-<!-- Remove if not applicable -->
-
-Fixes # (issue)
-
-
-## Before submitting
-- [ ] This PR fixes a typo or improves the docs (you can dismiss the other checks if that's the case).
-- [ ] Did you read the [contributor guideline](https://github.com/huggingface/diffusers/blob/main/CONTRIBUTING.md)?
-- [ ] Did you read our [philosophy doc](https://github.com/huggingface/diffusers/blob/main/PHILOSOPHY.md) (important for complex PRs)?
-- [ ] Was this discussed/approved via a Github issue or the [forum](https://discuss.huggingface.co/)? Please add a link to it if that's the case.
-- [ ] Did you make sure to update the documentation with your changes? Here are the
-      [documentation guidelines](https://github.com/huggingface/diffusers/tree/main/docs), and
-      [here are tips on formatting docstrings](https://github.com/huggingface/transformers/tree/main/docs#writing-source-documentation).
-- [ ] Did you write any new necessary tests?
-
-
-## Who can review?
-
-Anyone in the community is free to review the PR once the tests have passed. Feel free to tag
-members/contributors who may be interested in your PR.
-
-<!-- Your PR will be replied to more quickly if you can figure out the right person to tag with @
-
- If you know how to use git blame, that is the easiest way, otherwise, here is a rough guide of **who to tag**.
- Please tag fewer than 3 people.
-
-Core library:
-
-- Schedulers: @williamberman and @patrickvonplaten
-- Pipelines:  @patrickvonplaten and @sayakpaul
-- Training examples: @sayakpaul and @patrickvonplaten
-- Docs: @stevenliu and @yiyixu
-- JAX and MPS: @pcuenca
-- Audio: @sanchit-gandhi
-- General functionalities: @patrickvonplaten and @sayakpaul
-
-Integrations:
-
-- deepspeed: HF Trainer/Accelerate: @pacman100
-
-HF projects:
-
-- accelerate: [different repo](https://github.com/huggingface/accelerate)
-- datasets: [different repo](https://github.com/huggingface/datasets)
-- transformers: [different repo](https://github.com/huggingface/transformers)
-- safetensors: [different repo](https://github.com/huggingface/safetensors)
-
--->

.github/workflows/build_documentation.yml

@@ -5,19 +5,15 @@ on:
     branches:
       - main
       - doc-builder*
-      - v*-release
       - v*-patch
 
 jobs:
-  build:
+   build:
     uses: huggingface/doc-builder/.github/workflows/build_main_documentation.yml@main
     with:
       commit_sha: ${{ github.sha }}
-      install_libgl1: true
       package: diffusers
       notebook_folder: diffusers_doc
-      languages: en ko zh
-
+      languages: en ko
     secrets:
       token: ${{ secrets.HUGGINGFACE_PUSH }}
-      hf_token: ${{ secrets.HF_DOC_BUILD_PUSH }}

.github/workflows/build_pr_documentation.yml

@@ -9,10 +9,9 @@ concurrency:
 
 jobs:
   build:
-    uses: huggingface/doc-builder/.github/workflows/build_pr_documentation.yml@main
+    uses: huggingface/doc-builder/.github/workflows/build_pr_documentation.yml@@test_xenova_regex_optim
     with:
       commit_sha: ${{ github.event.pull_request.head.sha }}
       pr_number: ${{ github.event.number }}
-      install_libgl1: true
       package: diffusers
-      languages: en ko zh
+      languages: en ko

.github/workflows/delete_doc_comment.yml

@@ -1,14 +1,13 @@
-name: Delete doc comment
+name: Delete dev documentation
 
 on:
-  workflow_run:
-    workflows: ["Delete doc comment trigger"]
-    types:
-      - completed
+  pull_request:
+    types: [ closed ]
 
 
 jobs:
   delete:
     uses: huggingface/doc-builder/.github/workflows/delete_doc_comment.yml@main
-    secrets:
-      comment_bot_token: ${{ secrets.COMMENT_BOT_TOKEN }}
+    with:
+      pr_number: ${{ github.event.number }}
+      package: diffusers

.github/workflows/delete_doc_comment_trigger.yml

@@ -1,12 +0,0 @@
-name: Delete doc comment trigger
-
-on:
-  pull_request:
-    types: [ closed ]
-
-
-jobs:
-  delete:
-    uses: huggingface/doc-builder/.github/workflows/delete_doc_comment_trigger.yml@main
-    with:
-      pr_number: ${{ github.event.number }}

.github/workflows/pr_dependency_test.yml

@@ -1,32 +0,0 @@
-name: Run dependency tests
-
-on:
-  pull_request:
-    branches:
-      - main
-  push:
-    branches:
-      - main
-
-concurrency:
-  group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
-  cancel-in-progress: true
-
-jobs:
-  check_dependencies:
-    runs-on: ubuntu-latest
-    steps:
-      - uses: actions/checkout@v3
-      - name: Set up Python
-        uses: actions/setup-python@v4
-        with:
-          python-version: "3.7"
-      - name: Install dependencies
-        run: |
-          python -m pip install --upgrade pip
-          pip install -e .
-          pip install pytest
-      - name: Check for soft dependencies
-        run: |
-          pytest tests/others/test_dependencies.py
-      

.github/workflows/pr_tests.yml

@@ -62,7 +62,7 @@ jobs:
 
     - name: Install dependencies
       run: |
-        apt-get update && apt-get install libsndfile1-dev libgl1 -y
+        apt-get update && apt-get install libsndfile1-dev -y
         python -m pip install -e .[quality,test]
 
     - name: Environment
@@ -81,7 +81,7 @@ jobs:
       if: ${{ matrix.config.framework == 'pytorch_models' }}
       run: |
         python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-          -s -v -k "not Flax and not Onnx and not Dependency" \
+          -s -v -k "not Flax and not Onnx" \
           --make-reports=tests_${{ matrix.config.report }} \
           tests/models tests/schedulers tests/others
 

.github/workflows/push_tests.yml

@@ -17,7 +17,6 @@ jobs:
   run_slow_tests:
     strategy:
       fail-fast: false
-      max-parallel: 1
       matrix:
         config:
           - name: Slow PyTorch CUDA tests on Ubuntu
@@ -61,7 +60,6 @@ jobs:
 
     - name: Install dependencies
       run: |
-        apt-get update && apt-get install libsndfile1-dev libgl1 -y
         python -m pip install -e .[quality,test]
 
     - name: Environment

.github/workflows/push_tests_fast.yml

@@ -60,7 +60,7 @@ jobs:
 
     - name: Install dependencies
       run: |
-        apt-get update && apt-get install libsndfile1-dev libgl1 -y
+        apt-get update && apt-get install libsndfile1-dev -y
         python -m pip install -e .[quality,test]
 
     - name: Environment

.github/workflows/upload_pr_documentation.yml

@@ -1,16 +0,0 @@
-name: Upload PR Documentation
-
-on:
-  workflow_run:
-    workflows: ["Build PR Documentation"]
-    types:
-      - completed
-
-jobs:
-  build:
-    uses: huggingface/doc-builder/.github/workflows/upload_pr_documentation.yml@main
-    with:
-      package_name: diffusers
-    secrets:
-      hf_token: ${{ secrets.HF_DOC_BUILD_PUSH }}
-      comment_bot_token: ${{ secrets.COMMENT_BOT_TOKEN }}

CONTRIBUTING.md

@@ -125,14 +125,14 @@ Awesome! Tell us what problem it solved for you.
 
 You can open a feature request [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feature_request.md&title=).
 
-#### 2.3 Feedback.
+#### 2.3 Feedback. 
 
 Feedback about the library design and why it is good or not good helps the core maintainers immensely to build a user-friendly library. To understand the philosophy behind the current design philosophy, please have a look [here](https://huggingface.co/docs/diffusers/conceptual/philosophy). If you feel like a certain design choice does not fit with the current design philosophy, please explain why and how it should be changed. If a certain design choice follows the design philosophy too much, hence restricting use cases, explain why and how it should be changed.
 If a certain design choice is very useful for you, please also leave a note as this is great feedback for future design decisions.
 
 You can open an issue about feedback [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
 
-#### 2.4 Technical questions.
+#### 2.4 Technical questions. 
 
 Technical questions are mainly about why certain code of the library was written in a certain way, or what a certain part of the code does. Please make sure to link to the code in question and please provide detail on
 why this part of the code is difficult to understand.
@@ -394,8 +394,8 @@ passes. You should run the tests impacted by your changes like this:
  ```bash
  $ pytest tests/<TEST_TO_RUN>.py
  ```
-
-Before you run the tests, please make sure you install the dependencies required for testing. You can do so
+ 
+Before you run the tests, please make sure you install the dependencies required for testing. You can do so 
 with this command:
 
  ```bash

PHILOSOPHY.md

@@ -27,18 +27,18 @@ In a nutshell, Diffusers is built to be a natural extension of PyTorch. Therefor
 
 ## Simple over easy
 
-As PyTorch states, **explicit is better than implicit** and **simple is better than complex**. This design philosophy is reflected in multiple parts of the library:
+As PyTorch states, **explicit is better than implicit** and **simple is better than complex**. This design philosophy is reflected in multiple parts of the library: 
 - We follow PyTorch's API with methods like [`DiffusionPipeline.to`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.to) to let the user handle device management.
 - Raising concise error messages is preferred to silently correct erroneous input. Diffusers aims at teaching the user, rather than making the library as easy to use as possible.
 - Complex model vs. scheduler logic is exposed instead of magically handled inside. Schedulers/Samplers are separated from diffusion models with minimal dependencies on each other. This forces the user to write the unrolled denoising loop. However, the separation allows for easier debugging and gives the user more control over adapting the denoising process or switching out diffusion models or schedulers.
-- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the unet, and the variational autoencoder, each have their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. Dreambooth or textual inversion training
+- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the unet, and the variational autoencoder, each have their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. Dreambooth or textual inversion training 
 is very simple thanks to diffusers' ability to separate single components of the diffusion pipeline.
 
 ## Tweakable, contributor-friendly over abstraction
 
-For large parts of the library, Diffusers adopts an important design principle of the [Transformers library](https://github.com/huggingface/transformers), which is to prefer copy-pasted code over hasty abstractions. This design principle is very opinionated and stands in stark contrast to popular design principles such as [Don't repeat yourself (DRY)](https://en.wikipedia.org/wiki/Don%27t_repeat_yourself).
+For large parts of the library, Diffusers adopts an important design principle of the [Transformers library](https://github.com/huggingface/transformers), which is to prefer copy-pasted code over hasty abstractions. This design principle is very opinionated and stands in stark contrast to popular design principles such as [Don't repeat yourself (DRY)](https://en.wikipedia.org/wiki/Don%27t_repeat_yourself). 
 In short, just like Transformers does for modeling files, diffusers prefers to keep an extremely low level of abstraction and very self-contained code for pipelines and schedulers.
-Functions, long code blocks, and even classes can be copied across multiple files which at first can look like a bad, sloppy design choice that makes the library unmaintainable.
+Functions, long code blocks, and even classes can be copied across multiple files which at first can look like a bad, sloppy design choice that makes the library unmaintainable. 
 **However**, this design has proven to be extremely successful for Transformers and makes a lot of sense for community-driven, open-source machine learning libraries because:
 - Machine Learning is an extremely fast-moving field in which paradigms, model architectures, and algorithms are changing rapidly, which therefore makes it very difficult to define long-lasting code abstractions.
 - Machine Learning practitioners like to be able to quickly tweak existing code for ideation and research and therefore prefer self-contained code over one that contains many abstractions.
@@ -47,10 +47,10 @@ Functions, long code blocks, and even classes can be copied across multiple file
 At Hugging Face, we call this design the **single-file policy** which means that almost all of the code of a certain class should be written in a single, self-contained file. To read more about the philosophy, you can have a look
 at [this blog post](https://huggingface.co/blog/transformers-design-philosophy).
 
-In diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such
+In diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such 
 as [DDPM](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/ddpm), [Stable Diffusion](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/stable_diffusion/overview#stable-diffusion-pipelines), [UnCLIP (Dalle-2)](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/unclip#overview) and [Imagen](https://imagen.research.google/) all rely on the same diffusion model, the [UNet](https://huggingface.co/docs/diffusers/api/models#diffusers.UNet2DConditionModel).
 
-Great, now you should have generally understood why 🧨 Diffusers is designed the way it is 🤗.
+Great, now you should have generally understood why 🧨 Diffusers is designed the way it is 🤗. 
 We try to apply these design principles consistently across the library. Nevertheless, there are some minor exceptions to the philosophy or some unlucky design choices. If you have feedback regarding the design, we would ❤️  to hear it [directly on GitHub](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
 
 ## Design Philosophy in Details
@@ -89,7 +89,7 @@ The following design principles are followed:
 - Models should by default have the highest precision and lowest performance setting.
 - To integrate new model checkpoints whose general architecture can be classified as an architecture that already exists in Diffusers, the existing model architecture shall be adapted to make it work with the new checkpoint. One should only create a new file if the model architecture is fundamentally different.
 - Models should be designed to be easily extendable to future changes. This can be achieved by limiting public function arguments, configuration arguments, and "foreseeing" future changes, *e.g.* it is usually better to add `string` "...type" arguments that can easily be extended to new future types instead of boolean `is_..._type` arguments. Only the minimum amount of changes shall be made to existing architectures to make a new model checkpoint work.
-- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and
+- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and 
 readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py).
 
 ### Schedulers
@@ -97,9 +97,9 @@ readable longterm, such as [UNet blocks](https://github.com/huggingface/diffuser
 Schedulers are responsible to guide the denoising process for inference as well as to define a noise schedule for training. They are designed as individual classes with loadable configuration files and strongly follow the **single-file policy**.
 
 The following design principles are followed:
-- All schedulers are found in [`src/diffusers/schedulers`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
-- Schedulers are **not** allowed to import from large utils files and shall be kept very self-contained.
-- One scheduler python file corresponds to one scheduler algorithm (as might be defined in a paper).
+- All schedulers are found in [`src/diffusers/schedulers`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers). 
+- Schedulers are **not** allowed to import from large utils files and shall be kept very self-contained. 
+- One scheduler python file corresponds to one scheduler algorithm (as might be defined in a paper). 
 - If schedulers share similar functionalities, we can make use of the `#Copied from` mechanism.
 - Schedulers all inherit from `SchedulerMixin` and `ConfigMixin`.
 - Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./using-diffusers/schedulers.mdx).

README.md

@@ -25,12 +25,12 @@
 
 ## Installation
 
-We recommend installing 🤗 Diffusers in a virtual environment from PyPi or Conda. For more details about installing [PyTorch](https://pytorch.org/get-started/locally/) and [Flax](https://flax.readthedocs.io/en/latest/#installation), please refer to their official documentation.
+We recommend installing 🤗 Diffusers in a virtual environment from PyPi or Conda. For more details about installing [PyTorch](https://pytorch.org/get-started/locally/) and [Flax](https://flax.readthedocs.io/en/latest/installation.html), please refer to their official documentation.
 
 ### PyTorch
 
 With `pip` (official package):
-
+    
 ```bash
 pip install --upgrade diffusers[torch]
 ```
@@ -107,7 +107,7 @@ Check out the [Quickstart](https://huggingface.co/docs/diffusers/quicktour) to l
 | [Training](https://huggingface.co/docs/diffusers/training/overview) | Guides for how to train a diffusion model for different tasks with different training techniques.                                                                                               |
 ## Contribution
 
-We ❤️  contributions from the open-source community!
+We ❤️  contributions from the open-source community! 
 If you want to contribute to this library, please check out our [Contribution guide](https://github.com/huggingface/diffusers/blob/main/CONTRIBUTING.md).
 You can look out for [issues](https://github.com/huggingface/diffusers/issues) you'd like to tackle to contribute to the library.
 - See [Good first issues](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22good+first+issue%22) for general opportunities to contribute
@@ -128,7 +128,7 @@ just hang out ☕.
   </tr>
   <tr style="border-top: 2px solid black">
     <td>Unconditional Image Generation</td>
-    <td><a href="https://huggingface.co/docs/diffusers/api/pipelines/ddpm"> DDPM </a></td>
+    <td><a href="https://huggingface.co/docs/diffusers/api/pipelines/ddpm"> DDPM </a></td> 
     <td><a href="https://huggingface.co/google/ddpm-ema-church-256"> google/ddpm-ema-church-256 </a></td>
   </tr>
   <tr style="border-top: 2px solid black">
@@ -185,13 +185,13 @@ just hang out ☕.
 
 ## Popular libraries using 🧨 Diffusers
 
-- https://github.com/microsoft/TaskMatrix
-- https://github.com/invoke-ai/InvokeAI
-- https://github.com/apple/ml-stable-diffusion
-- https://github.com/Sanster/lama-cleaner
+- https://github.com/microsoft/TaskMatrix 
+- https://github.com/invoke-ai/InvokeAI 
+- https://github.com/apple/ml-stable-diffusion   
+- https://github.com/Sanster/lama-cleaner 
 - https://github.com/IDEA-Research/Grounded-Segment-Anything
-- https://github.com/ashawkey/stable-dreamfusion
-- https://github.com/deep-floyd/IF
+- https://github.com/ashawkey/stable-dreamfusion 
+- https://github.com/deep-floyd/IF  
 - https://github.com/bentoml/BentoML
 - https://github.com/bmaltais/kohya_ss
 - +3000 other amazing GitHub repositories 💪

docker/diffusers-pytorch-cpu/Dockerfile

@@ -14,7 +14,6 @@ RUN apt update && \
                    libsndfile1-dev \
                    python3.8 \
                    python3-pip \
-                   libgl1 \
                    python3.8-venv && \
     rm -rf /var/lib/apt/lists
 
@@ -28,7 +27,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
         torch \
         torchvision \
         torchaudio \
-        invisible_watermark \
         --extra-index-url https://download.pytorch.org/whl/cpu && \
     python3 -m pip install --no-cache-dir \
         accelerate \
@@ -42,4 +40,4 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
         tensorboard \
         transformers
 
-CMD ["/bin/bash"]
+CMD ["/bin/bash"]

docker/diffusers-pytorch-cuda/Dockerfile

@@ -12,7 +12,6 @@ RUN apt update && \
                    curl \
                    ca-certificates \
                    libsndfile1-dev \
-                   libgl1 \
                    python3.8 \
                    python3-pip \
                    python3.8-venv && \
@@ -27,8 +26,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
     python3 -m pip install --no-cache-dir \
         torch \
         torchvision \
-        torchaudio \
-        invisible_watermark && \
+        torchaudio && \
     python3 -m pip install --no-cache-dir \
         accelerate \
         datasets \

docs/source/_config.py

@@ -6,4 +6,4 @@ INSTALL_CONTENT = """
 # ! pip install git+https://github.com/huggingface/diffusers.git
 """
 
-notebook_first_cells = [{"type": "code", "content": INSTALL_CONTENT}]
+notebook_first_cells = [{"type": "code", "content": INSTALL_CONTENT}]

docs/source/en/_toctree.yml

@@ -50,8 +50,6 @@
       title: Distributed inference with multiple GPUs
     - local: using-diffusers/reusing_seeds
       title: Improve image quality with deterministic generation
-    - local: using-diffusers/control_brightness
-      title: Control image brightness
     - local: using-diffusers/reproducibility
       title: Create reproducible pipelines
     - local: using-diffusers/custom_pipeline_examples
@@ -132,6 +130,8 @@
   title: Conceptual Guides
 - sections:
   - sections:
+    - local: api/models
+      title: Models
     - local: api/attnprocessor
       title: Attention Processor
     - local: api/diffusion_pipeline
@@ -144,48 +144,16 @@
       title: Outputs
     - local: api/loaders
       title: Loaders
-    - local: api/utilities
-      title: Utilities
-    - local: api/image_processor
-      title: VAE Image Processor
     title: Main Classes
-  - sections:
-    - local: api/models/overview
-      title: Overview
-    - local: api/models/unet
-      title: UNet1DModel
-    - local: api/models/unet2d
-      title: UNet2DModel
-    - local: api/models/unet2d-cond
-      title: UNet2DConditionModel
-    - local: api/models/unet3d-cond
-      title: UNet3DConditionModel
-    - local: api/models/vq
-      title: VQModel
-    - local: api/models/autoencoderkl
-      title: AutoencoderKL
-    - local: api/models/transformer2d
-      title: Transformer2D
-    - local: api/models/transformer_temporal
-      title: Transformer Temporal
-    - local: api/models/prior_transformer
-      title: Prior Transformer
-    - local: api/models/controlnet
-      title: ControlNet
-    title: Models
   - sections:
     - local: api/pipelines/overview
       title: Overview
     - local: api/pipelines/alt_diffusion
       title: AltDiffusion
-    - local: api/pipelines/attend_and_excite
-      title: Attend and Excite
     - local: api/pipelines/audio_diffusion
       title: Audio Diffusion
     - local: api/pipelines/audioldm
       title: AudioLDM
-    - local: api/pipelines/consistency_models
-      title: Consistency Models
     - local: api/pipelines/controlnet
       title: ControlNet
     - local: api/pipelines/cycle_diffusion
@@ -196,38 +164,26 @@
       title: DDIM
     - local: api/pipelines/ddpm
       title: DDPM
-    - local: api/pipelines/diffedit
-      title: DiffEdit
     - local: api/pipelines/dit
       title: DiT
     - local: api/pipelines/if
       title: IF
-    - local: api/pipelines/pix2pix
-      title: InstructPix2Pix
     - local: api/pipelines/kandinsky
       title: Kandinsky
     - local: api/pipelines/latent_diffusion
       title: Latent Diffusion
-    - local: api/pipelines/panorama
-      title: MultiDiffusion Panorama
     - local: api/pipelines/paint_by_example
       title: PaintByExample
-    - local: api/pipelines/paradigms
-      title: Parallel Sampling of Diffusion Models
-    - local: api/pipelines/pix2pix_zero
-      title: Pix2Pix Zero
     - local: api/pipelines/pndm
       title: PNDM
     - local: api/pipelines/repaint
       title: RePaint
+    - local: api/pipelines/stable_diffusion_safe
+      title: Safe Stable Diffusion
     - local: api/pipelines/score_sde_ve
       title: Score SDE VE
-    - local: api/pipelines/self_attention_guidance
-      title: Self-Attention Guidance
     - local: api/pipelines/semantic_stable_diffusion
       title: Semantic Guidance
-    - local: api/pipelines/shap_e
-      title: Shap-E
     - local: api/pipelines/spectrogram_diffusion
       title: Spectrogram Diffusion
     - sections:
@@ -243,25 +199,31 @@
         title: Depth-to-Image
       - local: api/pipelines/stable_diffusion/image_variation
         title: Image-Variation
-      - local: api/pipelines/stable_diffusion/stable_diffusion_safe
-        title: Safe Stable Diffusion
-      - local: api/pipelines/stable_diffusion/stable_diffusion_2
-        title: Stable Diffusion 2
-      - local: api/pipelines/stable_diffusion/stable_diffusion_xl
-        title: Stable Diffusion XL
-      - local: api/pipelines/stable_diffusion/latent_upscale
-        title: Stable-Diffusion-Latent-Upscaler
       - local: api/pipelines/stable_diffusion/upscale
         title: Super-Resolution
-      - local: api/pipelines/stable_diffusion/ldm3d_diffusion
-        title: LDM3D Text-to-(RGB, Depth)
+      - local: api/pipelines/stable_diffusion/latent_upscale
+        title: Stable-Diffusion-Latent-Upscaler
+      - local: api/pipelines/stable_diffusion/pix2pix
+        title: InstructPix2Pix
+      - local: api/pipelines/stable_diffusion/attend_and_excite
+        title: Attend and Excite
+      - local: api/pipelines/stable_diffusion/pix2pix_zero
+        title: Pix2Pix Zero
+      - local: api/pipelines/stable_diffusion/self_attention_guidance
+        title: Self-Attention Guidance
+      - local: api/pipelines/stable_diffusion/panorama
+        title: MultiDiffusion Panorama
+      - local: api/pipelines/stable_diffusion/model_editing
+        title: Text-to-Image Model Editing
+      - local: api/pipelines/stable_diffusion/diffedit
+        title: DiffEdit
       title: Stable Diffusion
+    - local: api/pipelines/stable_diffusion_2
+      title: Stable Diffusion 2
     - local: api/pipelines/stable_unclip
       title: Stable unCLIP
     - local: api/pipelines/stochastic_karras_ve
       title: Stochastic Karras VE
-    - local: api/pipelines/model_editing
-      title: Text-to-Image Model Editing
     - local: api/pipelines/text_to_video
       title: Text-to-Video
     - local: api/pipelines/text_to_video_zero
@@ -280,8 +242,6 @@
   - sections:
     - local: api/schedulers/overview
       title: Overview
-    - local: api/schedulers/cm_stochastic_iterative
-      title: Consistency Model Multistep Scheduler
     - local: api/schedulers/ddim
       title: DDIM
     - local: api/schedulers/ddim_inverse

docs/source/en/api/attnprocessor.mdx

@@ -11,9 +11,6 @@ An attention processor is a class for applying different types of attention mech
 ## LoRAAttnProcessor
 [[autodoc]] models.attention_processor.LoRAAttnProcessor
 
-## LoRAAttnProcessor2_0
-[[autodoc]] models.attention_processor.LoRAAttnProcessor2_0
-
 ## CustomDiffusionAttnProcessor
 [[autodoc]] models.attention_processor.CustomDiffusionAttnProcessor
 

docs/source/en/api/configuration.mdx

@@ -12,13 +12,8 @@ specific language governing permissions and limitations under the License.
 
 # Configuration
 
-Schedulers from [`~schedulers.scheduling_utils.SchedulerMixin`] and models from [`ModelMixin`] inherit from [`ConfigMixin`] which stores all the parameters that are passed to their respective `__init__` methods in a JSON-configuration file.
-
-<Tip>
-
-To use private or [gated](https://huggingface.co/docs/hub/models-gated#gated-models) models, log-in with `huggingface-cli login`.
-
-</Tip>
+Schedulers from [`~schedulers.scheduling_utils.SchedulerMixin`] and models from [`ModelMixin`] inherit from [`ConfigMixin`] which conveniently takes care of storing all the parameters that are 
+passed to their respective `__init__` methods in a JSON-configuration file.
 
 ## ConfigMixin
 

docs/source/en/api/diffusion_pipeline.mdx

@@ -12,25 +12,41 @@ specific language governing permissions and limitations under the License.
 
 # Pipelines
 
-The [`DiffusionPipeline`] is the quickest way to load any pretrained diffusion pipeline from the [Hub](https://huggingface.co/models?library=diffusers) for inference.
+The [`DiffusionPipeline`] is the easiest way to load any pretrained diffusion pipeline from the [Hub](https://huggingface.co/models?library=diffusers) and to use it in inference.
 
 <Tip>
-
-You shouldn't use the [`DiffusionPipeline`] class for training or finetuning a diffusion model. Individual 
-components (for example, [`UNet2DModel`] and [`UNet2DConditionModel`]) of diffusion pipelines are usually trained individually, so we suggest directly working with them instead.
+	
+	One should not use the Diffusion Pipeline class for training or fine-tuning a diffusion model. Individual 
+	components of diffusion pipelines are usually trained individually, so we suggest to directly work 
+	with [`UNetModel`] and [`UNetConditionModel`].
 
 </Tip>
 
-The pipeline type (for example [`StableDiffusionPipeline`]) of any diffusion pipeline loaded with [`~DiffusionPipeline.from_pretrained`] is automatically 
-detected and pipeline components are loaded and passed to the `__init__` function of the pipeline.
+Any diffusion pipeline that is loaded with [`~DiffusionPipeline.from_pretrained`] will automatically 
+detect the pipeline type, *e.g.* [`StableDiffusionPipeline`] and consequently load each component of the 
+pipeline and pass them into the `__init__` function of the pipeline, *e.g.* [`~StableDiffusionPipeline.__init__`].
 
 Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrained`].
 
 ## DiffusionPipeline
-
 [[autodoc]] DiffusionPipeline
 	- all
 	- __call__
 	- device
 	- to
 	- components
+
+## ImagePipelineOutput
+By default diffusion pipelines return an object of class
+
+[[autodoc]] pipelines.ImagePipelineOutput
+
+## AudioPipelineOutput
+By default diffusion pipelines return an object of class
+
+[[autodoc]] pipelines.AudioPipelineOutput
+
+## ImageTextPipelineOutput
+By default diffusion pipelines return an object of class
+
+[[autodoc]] ImageTextPipelineOutput

docs/source/en/api/image_processor.mdx

@@ -1,27 +0,0 @@
-<!--Copyright 2023 The HuggingFace Team. All rights reserved.
-
-Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
-the License. You may obtain a copy of the License at
-
-http://www.apache.org/licenses/LICENSE-2.0
-
-Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
-an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
-specific language governing permissions and limitations under the License.
--->
-
-# VAE Image Processor
-
-The [`VaeImageProcessor`] provides a unified API for [`StableDiffusionPipeline`]'s to prepare image inputs for VAE encoding and post-processing outputs once they're decoded. This includes transformations such as resizing, normalization, and conversion between PIL Image, PyTorch, and NumPy arrays. 
-
-All pipelines with [`VaeImageProcessor`] accepts PIL Image, PyTorch tensor, or NumPy arrays as image inputs and returns outputs based on the `output_type` argument by the user. You can pass encoded image latents directly to the pipeline and return latents from the pipeline as a specific output with the `output_type` argument (for example `output_type="pt"`). This allows you to take the generated latents from one pipeline and pass it to another pipeline as input without leaving the latent space. It also makes it much easier to use multiple pipelines together by passing PyTorch tensors directly between different pipelines. 
-
-## VaeImageProcessor
-
-[[autodoc]] image_processor.VaeImageProcessor
-
-## VaeImageProcessorLDM3D
-
-The [`VaeImageProcessorLDM3D`] accepts RGB and depth inputs and returns RGB and depth outputs.
-
-[[autodoc]] image_processor.VaeImageProcessorLDM3D

docs/source/en/api/loaders.mdx

@@ -12,26 +12,31 @@ specific language governing permissions and limitations under the License.
 
 # Loaders
 
-Adapters (textual inversion, LoRA, hypernetworks) allow you to modify a diffusion model to generate images in a specific style without training or finetuning the entire model. The adapter weights are typically only a tiny fraction of the pretrained model's which making them very portable. 🤗 Diffusers provides an easy-to-use `LoaderMixin` API to load adapter weights.
+There are many ways to train adapter neural networks for diffusion models, such as 
+- [Textual Inversion](./training/text_inversion.mdx)
+- [LoRA](https://github.com/cloneofsimo/lora)
+- [Hypernetworks](https://arxiv.org/abs/1609.09106)
 
-<Tip warning={true}>
+Such adapter neural networks often only consist of a fraction of the number of weights compared 
+to the pretrained model and as such are very portable. The Diffusers library offers an easy-to-use
+API to load such adapter neural networks via the [`loaders.py` module](https://github.com/huggingface/diffusers/blob/main/src/diffusers/loaders.py). 
 
-🧪 The `LoaderMixins` are highly experimental and prone to future changes. To use private or [gated](https://huggingface.co/docs/hub/models-gated#gated-models) models, log-in with `huggingface-cli login`.
+**Note**: This module is still highly experimental and prone to future changes.
 
-</Tip>
+## LoaderMixins
 
-## UNet2DConditionLoadersMixin
+### UNet2DConditionLoadersMixin
 
 [[autodoc]] loaders.UNet2DConditionLoadersMixin
 
-## TextualInversionLoaderMixin
+### TextualInversionLoaderMixin
 
 [[autodoc]] loaders.TextualInversionLoaderMixin
 
-## LoraLoaderMixin
+### LoraLoaderMixin
 
 [[autodoc]] loaders.LoraLoaderMixin
 
-## FromSingleFileMixin
+### FromCkptMixin
 
-[[autodoc]] loaders.FromSingleFileMixin
+[[autodoc]] loaders.FromCkptMixin

docs/source/en/api/logging.mdx

@@ -12,9 +12,12 @@ specific language governing permissions and limitations under the License.
 
 # Logging
 
-🤗 Diffusers has a centralized logging system to easily manage the verbosity of the library. The default verbosity is set to `WARNING`.
+🧨 Diffusers has a centralized logging system, so that you can setup the verbosity of the library easily.
 
-To change the verbosity level, use one of the direct setters. For instance, to change the verbosity to the `INFO` level.
+Currently the default verbosity of the library is `WARNING`.
+
+To change the level of verbosity, just use one of the direct setters. For instance, here is how to change the verbosity
+to the INFO level.
 
 ```python
 import diffusers
@@ -30,7 +33,7 @@ DIFFUSERS_VERBOSITY=error ./myprogram.py
 ```
 
 Additionally, some `warnings` can be disabled by setting the environment variable
-`DIFFUSERS_NO_ADVISORY_WARNINGS` to a true value, like `1`. This disables any warning logged by
+`DIFFUSERS_NO_ADVISORY_WARNINGS` to a true value, like *1*. This will disable any warning that is logged using
 [`logger.warning_advice`]. For example:
 
 ```bash
@@ -49,21 +52,20 @@ logger.warning("WARN")
 ```
 
 
-All methods of the logging module are documented below. The main methods are
+All the methods of this logging module are documented below, the main ones are
 [`logging.get_verbosity`] to get the current level of verbosity in the logger and
-[`logging.set_verbosity`] to set the verbosity to the level of your choice. 
+[`logging.set_verbosity`] to set the verbosity to the level of your choice. In order (from the least
+verbose to the most verbose), those levels (with their corresponding int values in parenthesis) are:
 
-In order from the least verbose to the most verbose:
+- `diffusers.logging.CRITICAL` or `diffusers.logging.FATAL` (int value, 50): only report the most
+  critical errors.
+- `diffusers.logging.ERROR` (int value, 40): only report errors.
+- `diffusers.logging.WARNING` or `diffusers.logging.WARN` (int value, 30): only reports error and
+  warnings. This is the default level used by the library.
+- `diffusers.logging.INFO` (int value, 20): reports error, warnings and basic information.
+- `diffusers.logging.DEBUG` (int value, 10): report all information.
 
-|                                                    Method | Integer value |                                         Description |
-|----------------------------------------------------------:|--------------:|----------------------------------------------------:|
-| `diffusers.logging.CRITICAL` or `diffusers.logging.FATAL` |            50 |                only report the most critical errors |
-|                                 `diffusers.logging.ERROR` |            40 |                                  only report errors |
-|   `diffusers.logging.WARNING` or `diffusers.logging.WARN` |            30 |           only report errors and warnings (default) |
-|                                  `diffusers.logging.INFO` |            20 | only report errors, warnings, and basic information |
-|                                 `diffusers.logging.DEBUG` |            10 |                              report all information |
-
-By default, `tqdm` progress bars are displayed during model download. [`logging.disable_progress_bar`] and [`logging.enable_progress_bar`] are used to enable or disable this behavior.
+By default, `tqdm` progress bars will be displayed during model download. [`logging.disable_progress_bar`] and [`logging.enable_progress_bar`] can be used to suppress or unsuppress this behavior.
 
 ## Base setters
 

docs/source/en/api/models.mdx

@@ -0,0 +1,107 @@
+<!--Copyright 2023 The HuggingFace Team. All rights reserved.
+
+Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
+the License. You may obtain a copy of the License at
+
+http://www.apache.org/licenses/LICENSE-2.0
+
+Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
+an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
+specific language governing permissions and limitations under the License.
+-->
+
+# Models
+
+Diffusers contains pretrained models for popular algorithms and modules for creating the next set of diffusion models.
+The primary function of these models is to denoise an input sample, by modeling the distribution \\(p_{\theta}(x_{t-1}|x_{t})\\). 
+The models are built on the base class ['ModelMixin'] that is a `torch.nn.module` with basic functionality for saving and loading models both locally and from the HuggingFace hub.
+
+## ModelMixin
+[[autodoc]] ModelMixin
+
+## UNet2DOutput
+[[autodoc]] models.unet_2d.UNet2DOutput
+
+## UNet2DModel
+[[autodoc]] UNet2DModel
+
+## UNet1DOutput
+[[autodoc]] models.unet_1d.UNet1DOutput
+
+## UNet1DModel
+[[autodoc]] UNet1DModel
+
+## UNet2DConditionOutput
+[[autodoc]] models.unet_2d_condition.UNet2DConditionOutput
+
+## UNet2DConditionModel
+[[autodoc]] UNet2DConditionModel
+
+## UNet3DConditionOutput
+[[autodoc]] models.unet_3d_condition.UNet3DConditionOutput
+
+## UNet3DConditionModel
+[[autodoc]] UNet3DConditionModel
+
+## DecoderOutput
+[[autodoc]] models.vae.DecoderOutput
+
+## VQEncoderOutput
+[[autodoc]] models.vq_model.VQEncoderOutput
+
+## VQModel
+[[autodoc]] VQModel
+
+## AutoencoderKLOutput
+[[autodoc]] models.autoencoder_kl.AutoencoderKLOutput
+
+## AutoencoderKL
+[[autodoc]] AutoencoderKL
+
+## Transformer2DModel
+[[autodoc]] Transformer2DModel
+
+## Transformer2DModelOutput
+[[autodoc]] models.transformer_2d.Transformer2DModelOutput
+
+## TransformerTemporalModel
+[[autodoc]] models.transformer_temporal.TransformerTemporalModel
+
+## Transformer2DModelOutput
+[[autodoc]] models.transformer_temporal.TransformerTemporalModelOutput
+
+## PriorTransformer
+[[autodoc]] models.prior_transformer.PriorTransformer
+
+## PriorTransformerOutput
+[[autodoc]] models.prior_transformer.PriorTransformerOutput
+
+## ControlNetOutput
+[[autodoc]] models.controlnet.ControlNetOutput
+
+## ControlNetModel
+[[autodoc]] ControlNetModel
+
+## FlaxModelMixin
+[[autodoc]] FlaxModelMixin
+
+## FlaxUNet2DConditionOutput
+[[autodoc]] models.unet_2d_condition_flax.FlaxUNet2DConditionOutput
+
+## FlaxUNet2DConditionModel
+[[autodoc]] FlaxUNet2DConditionModel
+
+## FlaxDecoderOutput
+[[autodoc]] models.vae_flax.FlaxDecoderOutput
+
+## FlaxAutoencoderKLOutput
+[[autodoc]] models.vae_flax.FlaxAutoencoderKLOutput
+
+## FlaxAutoencoderKL
+[[autodoc]] FlaxAutoencoderKL
+
+## FlaxControlNetOutput
+[[autodoc]] models.controlnet_flax.FlaxControlNetOutput
+
+## FlaxControlNetModel
+[[autodoc]] FlaxControlNetModel

docs/source/en/api/models/autoencoderkl.mdx

@@ -1,31 +0,0 @@
-# AutoencoderKL
-
-The variational autoencoder (VAE) model with KL loss was introduced in [Auto-Encoding Variational Bayes](https://arxiv.org/abs/1312.6114v11) by Diederik P. Kingma and Max Welling. The model is used in 🤗 Diffusers to encode images into latents and to decode latent representations into images.
-
-The abstract from the paper is:
-
-*How can we perform efficient inference and learning in directed probabilistic models, in the presence of continuous latent variables with intractable posterior distributions, and large datasets? We introduce a stochastic variational inference and learning algorithm that scales to large datasets and, under some mild differentiability conditions, even works in the intractable case. Our contributions are two-fold. First, we show that a reparameterization of the variational lower bound yields a lower bound estimator that can be straightforwardly optimized using standard stochastic gradient methods. Second, we show that for i.i.d. datasets with continuous latent variables per datapoint, posterior inference can be made especially efficient by fitting an approximate inference model (also called a recognition model) to the intractable posterior using the proposed lower bound estimator. Theoretical advantages are reflected in experimental results.*
-
-## AutoencoderKL
-
-[[autodoc]] AutoencoderKL
-
-## AutoencoderKLOutput
-
-[[autodoc]] models.autoencoder_kl.AutoencoderKLOutput
-
-## DecoderOutput
-
-[[autodoc]] models.vae.DecoderOutput
-
-## FlaxAutoencoderKL
-
-[[autodoc]] FlaxAutoencoderKL
-
-## FlaxAutoencoderKLOutput
-
-[[autodoc]] models.vae_flax.FlaxAutoencoderKLOutput
-
-## FlaxDecoderOutput
-
-[[autodoc]] models.vae_flax.FlaxDecoderOutput

docs/source/en/api/models/controlnet.mdx

@@ -1,23 +0,0 @@
-# ControlNet
-
-The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
-
-The abstract from the paper is:
-
-*We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
-
-## ControlNetModel
-
-[[autodoc]] ControlNetModel
-
-## ControlNetOutput
-
-[[autodoc]] models.controlnet.ControlNetOutput
-
-## FlaxControlNetModel
-
-[[autodoc]] FlaxControlNetModel
-
-## FlaxControlNetOutput
-
-[[autodoc]] models.controlnet_flax.FlaxControlNetOutput

docs/source/en/api/models/overview.mdx

@@ -1,12 +0,0 @@
-# Models
-
-🤗 Diffusers provides pretrained models for popular algorithms and modules to create custom diffusion systems. The primary function of models is to denoise an input sample as modeled by the distribution \\(p_{\theta}(x_{t-1}|x_{t})\\).
-
-All models are built from the base [`ModelMixin`] class which is a [`torch.nn.module`](https://pytorch.org/docs/stable/generated/torch.nn.Module.html) providing basic functionality for saving and loading models, locally and from the Hugging Face Hub.
-
-## ModelMixin
-[[autodoc]] ModelMixin
-
-## FlaxModelMixin
-
-[[autodoc]] FlaxModelMixin

docs/source/en/api/models/prior_transformer.mdx

@@ -1,16 +0,0 @@
-# Prior Transformer
-
-The Prior Transformer was originally introduced in [Hierarchical Text-Conditional Image Generation with CLIP Latents
-](https://huggingface.co/papers/2204.06125) by Ramesh et al. It is used to predict CLIP image embeddings from CLIP text embeddings; image embeddings are predicted through a denoising diffusion process.
-
-The abstract from the paper is:
-
-*Contrastive models like CLIP have been shown to learn robust representations of images that capture both semantics and style. To leverage these representations for image generation, we propose a two-stage model: a prior that generates a CLIP image embedding given a text caption, and a decoder that generates an image conditioned on the image embedding. We show that explicitly generating image representations improves image diversity with minimal loss in photorealism and caption similarity. Our decoders conditioned on image representations can also produce variations of an image that preserve both its semantics and style, while varying the non-essential details absent from the image representation. Moreover, the joint embedding space of CLIP enables language-guided image manipulations in a zero-shot fashion. We use diffusion models for the decoder and experiment with both autoregressive and diffusion models for the prior, finding that the latter are computationally more efficient and produce higher-quality samples.*
-
-## PriorTransformer
-
-[[autodoc]] PriorTransformer
-
-## PriorTransformerOutput
-
-[[autodoc]] models.prior_transformer.PriorTransformerOutput

docs/source/en/api/models/transformer2d.mdx

@@ -1,29 +0,0 @@
-# Transformer2D
-
-A Transformer model for image-like data from [CompVis](https://huggingface.co/CompVis) that is based on the [Vision Transformer](https://huggingface.co/papers/2010.11929) introduced by Dosovitskiy et al. The [`Transformer2DModel`] accepts discrete (classes of vector embeddings) or continuous (actual embeddings) inputs.
-
-When the input is **continuous**:
-
-1. Project the input and reshape it to `(batch_size, sequence_length, feature_dimension)`.
-2. Apply the Transformer blocks in the standard way.
-3. Reshape to image.
-
-When the input is **discrete**:
-
-<Tip>
-
-It is assumed one of the input classes is the masked latent pixel. The predicted classes of the unnoised image don't contain a prediction for the masked pixel because the unnoised image cannot be masked.
-
-</Tip>
-
-1. Convert input (classes of latent pixels) to embeddings and apply positional embeddings.
-2. Apply the Transformer blocks in the standard way.
-3. Predict classes of unnoised image.
-
-## Transformer2DModel
-
-[[autodoc]] Transformer2DModel
-
-## Transformer2DModelOutput
-
-[[autodoc]] models.transformer_2d.Transformer2DModelOutput

docs/source/en/api/models/transformer_temporal.mdx

@@ -1,11 +0,0 @@
-# Transformer Temporal
-
-A Transformer model for video-like data.
-
-## TransformerTemporalModel
-
-[[autodoc]] models.transformer_temporal.TransformerTemporalModel
-
-## TransformerTemporalModelOutput
-
-[[autodoc]] models.transformer_temporal.TransformerTemporalModelOutput

docs/source/en/api/models/unet.mdx

@@ -1,13 +0,0 @@
-# UNet1DModel
-
-The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 1D UNet model.
-
-The abstract from the paper is:
-
-*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
-
-## UNet1DModel
-[[autodoc]] UNet1DModel
-
-## UNet1DOutput
-[[autodoc]] models.unet_1d.UNet1DOutput

docs/source/en/api/models/unet2d-cond.mdx

@@ -1,19 +0,0 @@
-# UNet2DConditionModel
-
-The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet conditional model.
-
-The abstract from the paper is:
-
-*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
-
-## UNet2DConditionModel
-[[autodoc]] UNet2DConditionModel
-
-## UNet2DConditionOutput
-[[autodoc]] models.unet_2d_condition.UNet2DConditionOutput
-
-## FlaxUNet2DConditionModel
-[[autodoc]] models.unet_2d_condition_flax.FlaxUNet2DConditionModel
-
-## FlaxUNet2DConditionOutput
-[[autodoc]] models.unet_2d_condition_flax.FlaxUNet2DConditionOutput

docs/source/en/api/models/unet2d.mdx

@@ -1,13 +0,0 @@
-# UNet2DModel
-
-The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet model.
-
-The abstract from the paper is:
-
-*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
-
-## UNet2DModel
-[[autodoc]] UNet2DModel
-
-## UNet2DOutput
-[[autodoc]] models.unet_2d.UNet2DOutput

docs/source/en/api/models/unet3d-cond.mdx

@@ -1,13 +0,0 @@
-# UNet3DConditionModel
-
-The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 3D UNet conditional model.
-
-The abstract from the paper is:
-
-*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
-
-## UNet3DConditionModel
-[[autodoc]] UNet3DConditionModel
-
-## UNet3DConditionOutput
-[[autodoc]] models.unet_3d_condition.UNet3DConditionOutput

docs/source/en/api/models/vq.mdx

@@ -1,15 +0,0 @@
-# VQModel
-
-The VQ-VAE model was introduced in [Neural Discrete Representation Learning](https://huggingface.co/papers/1711.00937) by Aaron van den Oord, Oriol Vinyals and Koray Kavukcuoglu. The model is used in 🤗 Diffusers to decode latent representations into images. Unlike [`AutoencoderKL`], the [`VQModel`] works in a quantized latent space.
-
-The abstract from the paper is:
-
-*Learning useful representations without supervision remains a key challenge in machine learning. In this paper, we propose a simple yet powerful generative model that learns such discrete representations. Our model, the Vector Quantised-Variational AutoEncoder (VQ-VAE), differs from VAEs in two key ways: the encoder network outputs discrete, rather than continuous, codes; and the prior is learnt rather than static. In order to learn a discrete latent representation, we incorporate ideas from vector quantisation (VQ). Using the VQ method allows the model to circumvent issues of "posterior collapse" -- where the latents are ignored when they are paired with a powerful autoregressive decoder -- typically observed in the VAE framework. Pairing these representations with an autoregressive prior, the model can generate high quality images, videos, and speech as well as doing high quality speaker conversion and unsupervised learning of phonemes, providing further evidence of the utility of the learnt representations.*
-
-## VQModel
-
-[[autodoc]] VQModel
-
-## VQEncoderOutput
-
-[[autodoc]] models.vq_model.VQEncoderOutput

docs/source/en/api/outputs.mdx

@@ -10,11 +10,13 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
 
-# Outputs
+# BaseOutputs
 
-All models outputs are subclasses of [`~utils.BaseOutput`], data structures containing all the information returned by the model. The outputs can also be used as tuples or dictionaries.
+All models have outputs that are instances of subclasses of [`~utils.BaseOutput`]. Those are
+data structures containing all the information returned by the model, but that can also be used as tuples or
+dictionaries.
 
-For example:
+Let's see how this looks in an example:
 
 ```python
 from diffusers import DDIMPipeline
@@ -23,45 +25,31 @@ pipeline = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32")
 outputs = pipeline()
 ```
 
-The `outputs` object is a [`~pipelines.ImagePipelineOutput`] which means it has an image attribute.
+The `outputs` object is a [`~pipelines.ImagePipelineOutput`], as we can see in the
+documentation of that class below, it means it has an image attribute.
 
-You can access each attribute as you normally would or with a keyword lookup, and if that attribute is not returned by the model, you will get `None`:
+You can access each attribute as you would usually do, and if that attribute has not been returned by the model, you will get `None`:
 
 ```python
 outputs.images
+```
+
+or via keyword lookup
+
+```python
 outputs["images"]
 ```
 
-When considering the `outputs` object as a tuple, it only considers the attributes that don't have `None` values.
-For instance, retrieving an image by indexing into it returns the tuple `(outputs.images)`:
+When considering our `outputs` object as tuple, it only considers the attributes that don't have `None` values.
+Here for instance, we could retrieve images via indexing:
 
 ```python
 outputs[:1]
 ```
 
-<Tip>
-
-To check a specific pipeline or model output, refer to its corresponding API documentation.
-
-</Tip>
+which will return the tuple `(outputs.images)` for instance.
 
 ## BaseOutput
 
 [[autodoc]] utils.BaseOutput
     - to_tuple
-
-## ImagePipelineOutput
-
-[[autodoc]] pipelines.ImagePipelineOutput
-
-## FlaxImagePipelineOutput
-
-[[autodoc]] pipelines.pipeline_flax_utils.FlaxImagePipelineOutput
-
-## AudioPipelineOutput
-
-[[autodoc]] pipelines.AudioPipelineOutput
-
-## ImageTextPipelineOutput
-
-[[autodoc]] ImageTextPipelineOutput

docs/source/en/api/pipelines/audio_diffusion.mdx

@@ -43,7 +43,7 @@ pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-256").to(devic
 
 output = pipe()
 display(output.images[0])
-display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
+display(Audio(output.audios[0], rate=mel.get_sample_rate()))
 ```
 
 ### Latent Audio Diffusion

docs/source/en/api/pipelines/consistency_models.mdx

@@ -1,87 +0,0 @@
-# Consistency Models
-
-Consistency Models were proposed in [Consistency Models](https://arxiv.org/abs/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
-
-The abstract of the [paper](https://arxiv.org/pdf/2303.01469.pdf) is as follows:
-
-*Diffusion models have significantly advanced the fields of image, audio, and video generation, but they depend on an iterative sampling process that causes slow generation. To overcome this limitation, we propose consistency models, a new family of models that generate high quality samples by directly mapping noise to data. They support fast one-step generation by design, while still allowing multistep sampling to trade compute for sample quality. They also support zero-shot data editing, such as image inpainting, colorization, and super-resolution, without requiring explicit training on these tasks. Consistency models can be trained either by distilling pre-trained diffusion models, or as standalone generative models altogether. Through extensive experiments, we demonstrate that they outperform existing distillation techniques for diffusion models in one- and few-step sampling, achieving the new state-of-the-art FID of 3.55 on CIFAR-10 and 6.20 on ImageNet 64x64 for one-step generation. When trained in isolation, consistency models become a new family of generative models that can outperform existing one-step, non-adversarial generative models on standard benchmarks such as CIFAR-10, ImageNet 64x64 and LSUN 256x256. *
-
-Resources:
-
-* [Paper](https://arxiv.org/abs/2303.01469)
-* [Original Code](https://github.com/openai/consistency_models)
-
-Available Checkpoints are:
-- *cd_imagenet64_l2 (64x64 resolution)* [openai/consistency-model-pipelines](https://huggingface.co/openai/diffusers-cd_imagenet64_l2)
-- *cd_imagenet64_lpips (64x64 resolution)* [openai/diffusers-cd_imagenet64_lpips](https://huggingface.co/openai/diffusers-cd_imagenet64_lpips)
-- *ct_imagenet64 (64x64 resolution)* [openai/diffusers-ct_imagenet64](https://huggingface.co/openai/diffusers-ct_imagenet64)
-- *cd_bedroom256_l2 (256x256 resolution)* [openai/diffusers-cd_bedroom256_l2](https://huggingface.co/openai/diffusers-cd_bedroom256_l2)
-- *cd_bedroom256_lpips (256x256 resolution)* [openai/diffusers-cd_bedroom256_lpips](https://huggingface.co/openai/diffusers-cd_bedroom256_lpips)
-- *ct_bedroom256 (256x256 resolution)* [openai/diffusers-ct_bedroom256](https://huggingface.co/openai/diffusers-ct_bedroom256)
-- *cd_cat256_l2 (256x256 resolution)* [openai/diffusers-cd_cat256_l2](https://huggingface.co/openai/diffusers-cd_cat256_l2)
-- *cd_cat256_lpips (256x256 resolution)* [openai/diffusers-cd_cat256_lpips](https://huggingface.co/openai/diffusers-cd_cat256_lpips)
-- *ct_cat256 (256x256 resolution)* [openai/diffusers-ct_cat256](https://huggingface.co/openai/diffusers-ct_cat256)
-
-## Available Pipelines
-
-| Pipeline | Tasks | Demo | Colab |
-|:---:|:---:|:---:|:---:|
-| [ConsistencyModelPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/pipeline_consistency_models.py) | *Unconditional Image Generation* | | |
-
-This pipeline was contributed by our community members [dg845](https://github.com/dg845) and [ayushtues](https://huggingface.co/ayushtues) ❤️
-
-## Usage Example
-
-```python
-import torch
-
-from diffusers import ConsistencyModelPipeline
-
-device = "cuda"
-# Load the cd_imagenet64_l2 checkpoint.
-model_id_or_path = "openai/diffusers-cd_imagenet64_l2"
-pipe = ConsistencyModelPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
-pipe.to(device)
-
-# Onestep Sampling
-image = pipe(num_inference_steps=1).images[0]
-image.save("consistency_model_onestep_sample.png")
-
-# Onestep sampling, class-conditional image generation
-# ImageNet-64 class label 145 corresponds to king penguins
-image = pipe(num_inference_steps=1, class_labels=145).images[0]
-image.save("consistency_model_onestep_sample_penguin.png")
-
-# Multistep sampling, class-conditional image generation
-# Timesteps can be explicitly specified; the particular timesteps below are from the original Github repo.
-# https://github.com/openai/consistency_models/blob/main/scripts/launch.sh#L77
-image = pipe(timesteps=[22, 0], class_labels=145).images[0]
-image.save("consistency_model_multistep_sample_penguin.png")
-```
-
-For an additional speed-up, one can also make use of `torch.compile`. Multiple images can be generated in <1 second as follows:
-
-```py
-import torch
-from diffusers import ConsistencyModelPipeline
-
-device = "cuda"
-# Load the cd_bedroom256_lpips checkpoint.
-model_id_or_path = "openai/diffusers-cd_bedroom256_lpips"
-pipe = ConsistencyModelPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
-pipe.to(device)
-
-pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
-
-# Multistep sampling
-# Timesteps can be explicitly specified; the particular timesteps below are from the original Github repo:
-# https://github.com/openai/consistency_models/blob/main/scripts/launch.sh#L83
-for _ in range(10):
-    image = pipe(timesteps=[17, 0]).images[0]
-    image.show()
-```
-
-## ConsistencyModelPipeline
-[[autodoc]] ConsistencyModelPipeline
-    - all
-    - __call__

docs/source/en/api/pipelines/kandinsky.mdx

@@ -11,94 +11,89 @@ specific language governing permissions and limitations under the License.
 
 ## Overview
 
-Kandinsky inherits best practices from [DALL-E 2](https://huggingface.co/papers/2204.06125) and [Latent Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/latent_diffusion), while introducing some new ideas.
+Kandinsky 2.1 inherits best practices from [DALL-E 2](https://arxiv.org/abs/2204.06125) and [Latent Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/latent_diffusion), while introducing some new ideas.
 
 It uses [CLIP](https://huggingface.co/docs/transformers/model_doc/clip) for encoding images and text, and a diffusion image prior (mapping) between latent spaces of CLIP modalities. This approach enhances the visual performance of the model and unveils new horizons in blending images and text-guided image manipulation.
 
-The Kandinsky model is created by [Arseniy Shakhmatov](https://github.com/cene555), [Anton Razzhigaev](https://github.com/razzant), [Aleksandr Nikolich](https://github.com/AlexWortega), [Igor Pavlov](https://github.com/boomb0om), [Andrey Kuznetsov](https://github.com/kuznetsoffandrey) and [Denis Dimitrov](https://github.com/denndimitrov). The original codebase can be found [here](https://github.com/ai-forever/Kandinsky-2)
+The Kandinsky model is created by [Arseniy Shakhmatov](https://github.com/cene555), [Anton Razzhigaev](https://github.com/razzant), [Aleksandr Nikolich](https://github.com/AlexWortega), [Igor Pavlov](https://github.com/boomb0om), [Andrey Kuznetsov](https://github.com/kuznetsoffandrey) and [Denis Dimitrov](https://github.com/denndimitrov) and the original codebase can be found [here](https://github.com/ai-forever/Kandinsky-2)
 
+## Available Pipelines:
+
+| Pipeline | Tasks | Colab
+|---|---|:---:|
+| [pipeline_kandinsky.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky.py) | *Text-to-Image Generation* | - |
+| [pipeline_kandinsky_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_inpaint.py) | *Image-Guided Image Generation* | - |
+| [pipeline_kandinsky_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_img2img.py) | *Image-Guided Image Generation* | - |
 
 ## Usage example
 
-In the following, we will walk you through some examples of how to use the Kandinsky pipelines to create some visually aesthetic artwork.
+In the following, we will walk you through some cool examples of using the Kandinsky pipelines to create some visually aesthetic artwork.
 
 ### Text-to-Image Generation
 
-For text-to-image generation, we need to use both [`KandinskyPriorPipeline`] and [`KandinskyPipeline`].
-The first step is to encode text prompts with CLIP and then diffuse the CLIP text embeddings to CLIP image embeddings,
-as first proposed in [DALL-E 2](https://cdn.openai.com/papers/dall-e-2.pdf).
-Let's throw a fun prompt at Kandinsky to see what it comes up with.
+For text-to-image generation, we need to use both [`KandinskyPriorPipeline`] and [`KandinskyPipeline`]. The first step is to encode text prompts with CLIP and then diffuse the CLIP text embeddings to CLIP image embeddings, as first proposed in [DALL-E 2](https://cdn.openai.com/papers/dall-e-2.pdf). Let's throw a fun prompt at Kandinsky to see what it comes up with :)
 
-```py
-prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
-```
-
-First, let's instantiate the prior pipeline and the text-to-image pipeline. Both 
-pipelines are diffusion models.
-
-
-```py
-from diffusers import DiffusionPipeline
-import torch
-
-pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16)
-pipe_prior.to("cuda")
-
-t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
-t2i_pipe.to("cuda")
-```
-
-<Tip warning={true}>
-
-By default, the text-to-image pipeline use [`DDIMScheduler`], you can change the scheduler to [`DDPMScheduler`]
-
-```py
-scheduler = DDPMScheduler.from_pretrained("kandinsky-community/kandinsky-2-1", subfolder="ddpm_scheduler")
-t2i_pipe = DiffusionPipeline.from_pretrained(
-    "kandinsky-community/kandinsky-2-1", scheduler=scheduler, torch_dtype=torch.float16
-)
-t2i_pipe.to("cuda")
-```
-
-</Tip>
-
-Now we pass the prompt through the prior to generate image embeddings. The prior
-returns both the image embeddings corresponding to the prompt and negative/unconditional image 
-embeddings corresponding to an empty string.
-
-```py
-image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
-```
-
-<Tip warning={true}>
-
-The text-to-image pipeline expects both `image_embeds`, `negative_image_embeds` and the original 
-`prompt` as the text-to-image pipeline uses another text encoder to better guide the second diffusion 
-process of `t2i_pipe`.
-
-By default, the prior returns unconditioned negative image embeddings corresponding to the negative prompt of `""`.
-For better results, you can also pass a `negative_prompt` to the prior. This will increase the effective batch size
-of the prior by a factor of 2.
-
-```py
+```python
 prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
 negative_prompt = "low quality, bad quality"
-
-image_embeds, negative_image_embeds = pipe_prior(prompt, negative_prompt, guidance_scale=1.0).to_tuple()
 ```
 
-</Tip>
+We will pass both the `prompt` and `negative_prompt` to our prior diffusion pipeline. In contrast to other diffusion pipelines, such as Stable Diffusion, the `prompt` and `negative_prompt` shall be passed separately so that we can retrieve a CLIP image embedding for each prompt input. You can use `guidance_scale`, and `num_inference_steps` arguments to guide this process, just like how you would normally do with all other pipelines in diffusers.
+
+```python
+from diffusers import KandinskyPriorPipeline
+import torch
+
+# create prior
+pipe_prior = KandinskyPriorPipeline.from_pretrained(
+    "kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16
+)
+pipe_prior.to("cuda")
+
+generator = torch.Generator(device="cuda").manual_seed(12)
+image_emb = pipe_prior(
+    prompt, guidance_scale=1.0, num_inference_steps=25, generator=generator, negative_prompt=negative_prompt
+).images
+
+zero_image_emb = pipe_prior(
+    negative_prompt, guidance_scale=1.0, num_inference_steps=25, generator=generator, negative_prompt=negative_prompt
+).images
+```
+
+Once we create the image embedding, we can use [`KandinskyPipeline`] to generate images.
+
+```python
+from PIL import Image
+from diffusers import KandinskyPipeline
 
 
-Next, we can pass the embeddings as well as the prompt to the text-to-image pipeline. Remember that 
-in case you are using a customized negative prompt, that you should pass this one also to the text-to-image pipelines
-with `negative_prompt=negative_prompt`:
+def image_grid(imgs, rows, cols):
+    assert len(imgs) == rows * cols
 
-```py
-image = t2i_pipe(
-    prompt, image_embeds=image_embeds, negative_image_embeds=negative_image_embeds, height=768, width=768
-).images[0]
-image.save("cheeseburger_monster.png")
+    w, h = imgs[0].size
+    grid = Image.new("RGB", size=(cols * w, rows * h))
+    grid_w, grid_h = grid.size
+
+    for i, img in enumerate(imgs):
+        grid.paste(img, box=(i % cols * w, i // cols * h))
+    return grid
+
+
+# create diffuser pipeline
+pipe = KandinskyPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
+pipe.to("cuda")
+
+images = pipe(
+    prompt,
+    image_embeds=image_emb,
+    negative_image_embeds=zero_image_emb,
+    num_images_per_prompt=2,
+    height=768,
+    width=768,
+    num_inference_steps=100,
+    guidance_scale=4.0,
+    generator=generator,
+).images
 ```
 
 One cheeseburger monster coming up! Enjoy! 
@@ -128,7 +123,6 @@ prompt = "birds eye view of a quilted paper style alien planet landscape, vibran
 ![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/alienplanet.png)
 
 
-
 ### Text Guided Image-to-Image Generation
 
 The same Kandinsky model weights can be used for text-guided image-to-image translation. In this case, just make sure to load the weights using the [`KandinskyImg2ImgPipeline`] pipeline.
@@ -169,15 +163,23 @@ pipe.to("cuda")
 prompt = "A fantasy landscape, Cinematic lighting"
 negative_prompt = "low quality, bad quality"
 
-image_embeds, negative_image_embeds = pipe_prior(prompt, negative_prompt).to_tuple()
+generator = torch.Generator(device="cuda").manual_seed(30)
+image_emb = pipe_prior(
+    prompt, guidance_scale=4.0, num_inference_steps=25, generator=generator, negative_prompt=negative_prompt
+).images
+
+zero_image_emb = pipe_prior(
+    negative_prompt, guidance_scale=4.0, num_inference_steps=25, generator=generator, negative_prompt=negative_prompt
+).images
 
 out = pipe(
     prompt,
     image=original_image,
-    image_embeds=image_embeds,
-    negative_image_embeds=negative_image_embeds,
+    image_embeds=image_emb,
+    negative_image_embeds=zero_image_emb,
     height=768,
     width=768,
+    num_inference_steps=500,
     strength=0.3,
 )
 
@@ -191,7 +193,7 @@ out.images[0].save("fantasy_land.png")
 
 You can use [`KandinskyInpaintPipeline`] to edit images. In this example, we will add a hat to the portrait of a cat.
 
-```py
+```python
 from diffusers import KandinskyInpaintPipeline, KandinskyPriorPipeline
 from diffusers.utils import load_image
 import torch
@@ -203,7 +205,7 @@ pipe_prior = KandinskyPriorPipeline.from_pretrained(
 pipe_prior.to("cuda")
 
 prompt = "a hat"
-prior_output = pipe_prior(prompt)
+image_emb, zero_image_emb = pipe_prior(prompt, return_dict=False)
 
 pipe = KandinskyInpaintPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-inpaint", torch_dtype=torch.float16)
 pipe.to("cuda")
@@ -220,7 +222,8 @@ out = pipe(
     prompt,
     image=init_image,
     mask_image=mask,
-    **prior_output,
+    image_embeds=image_emb,
+    negative_image_embeds=zero_image_emb,
     height=768,
     width=768,
     num_inference_steps=150,
@@ -243,6 +246,7 @@ from diffusers.utils import load_image
 import PIL
 
 import torch
+from torchvision import transforms
 
 pipe_prior = KandinskyPriorPipeline.from_pretrained(
     "kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16
@@ -259,356 +263,44 @@ img2 = load_image(
 
 # add all the conditions we want to interpolate, can be either text or image
 images_texts = ["a cat", img1, img2]
-
 # specify the weights for each condition in images_texts
 weights = [0.3, 0.3, 0.4]
-
-# We can leave the prompt empty
-prompt = ""
-prior_out = pipe_prior.interpolate(images_texts, weights)
+image_emb, zero_image_emb = pipe_prior.interpolate(images_texts, weights)
 
 pipe = KandinskyPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
 pipe.to("cuda")
 
-image = pipe(prompt, **prior_out, height=768, width=768).images[0]
+image = pipe(
+    "", image_embeds=image_emb, negative_image_embeds=zero_image_emb, height=768, width=768, num_inference_steps=150
+).images[0]
 
 image.save("starry_cat.png")
 ```
 ![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/starry_cat.png)
 
 
-### Text-to-Image Generation with ControlNet Conditioning
-
-In the following, we give a simple example of how to use [`KandinskyV22ControlnetPipeline`] to add control to the text-to-image generation with a depth image.
-
-First, let's take an image and extract its depth map.
-
-```python
-from diffusers.utils import load_image
-
-img = load_image(
-    "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/cat.png"
-).resize((768, 768))
-```
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/cat.png)
-
-We can use the `depth-estimation` pipeline from transformers to process the image and retrieve its depth map.
-
-```python
-import torch
-import numpy as np
-
-from transformers import pipeline
-from diffusers.utils import load_image
-
-
-def make_hint(image, depth_estimator):
-    image = depth_estimator(image)["depth"]
-    image = np.array(image)
-    image = image[:, :, None]
-    image = np.concatenate([image, image, image], axis=2)
-    detected_map = torch.from_numpy(image).float() / 255.0
-    hint = detected_map.permute(2, 0, 1)
-    return hint
-
-
-depth_estimator = pipeline("depth-estimation")
-hint = make_hint(img, depth_estimator).unsqueeze(0).half().to("cuda")
-```
-Now, we load the prior pipeline and the text-to-image controlnet pipeline
-
-```python
-from diffusers import KandinskyV22PriorPipeline, KandinskyV22ControlnetPipeline
-
-pipe_prior = KandinskyV22PriorPipeline.from_pretrained(
-    "kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16
-)
-pipe_prior = pipe_prior.to("cuda")
-
-pipe = KandinskyV22ControlnetPipeline.from_pretrained(
-    "kandinsky-community/kandinsky-2-2-controlnet-depth", torch_dtype=torch.float16
-)
-pipe = pipe.to("cuda")
-```
-
-We pass the prompt and negative prompt through the prior to generate image embeddings
-
-```python
-prompt = "A robot, 4k photo"
-
-negative_prior_prompt = "lowres, text, error, cropped, worst quality, low quality, jpeg artifacts, ugly, duplicate, morbid, mutilated, out of frame, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, mutation, deformed, blurry, dehydrated, bad anatomy, bad proportions, extra limbs, cloned face, disfigured, gross proportions, malformed limbs, missing arms, missing legs, extra arms, extra legs, fused fingers, too many fingers, long neck, username, watermark, signature"
-
-generator = torch.Generator(device="cuda").manual_seed(43)
-image_emb, zero_image_emb = pipe_prior(
-    prompt=prompt, negative_prompt=negative_prior_prompt, generator=generator
-).to_tuple()
-```
-
-Now we can pass the image embeddings and the depth image we extracted to the controlnet pipeline. With Kandinsky 2.2, only prior pipelines accept `prompt` input. You do not need to pass the prompt to the controlnet pipeline.
-
-```python
-images = pipe(
-    image_embeds=image_emb,
-    negative_image_embeds=zero_image_emb,
-    hint=hint,
-    num_inference_steps=50,
-    generator=generator,
-    height=768,
-    width=768,
-).images
-
-images[0].save("robot_cat.png")
-```
-
-The output image looks as follow:
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/robot_cat_text2img.png)
-
-### Image-to-Image Generation with ControlNet Conditioning
-
-Kandinsky 2.2 also includes a [`KandinskyV22ControlnetImg2ImgPipeline`] that will allow you to add control to the image generation process with both the image and its depth map. This pipeline works really well with [`KandinskyV22PriorEmb2EmbPipeline`], which generates image embeddings based on both a text prompt and an image. 
-
-For our robot cat example, we will pass the prompt and cat image together to the prior pipeline to generate an image embedding. We will then use that image embedding and the depth map of the cat to further control the image generation process. 
-
-We can use the same cat image and its depth map from the last example.
-
-```python
-import torch
-import numpy as np
-
-from diffusers import KandinskyV22PriorEmb2EmbPipeline, KandinskyV22ControlnetImg2ImgPipeline
-from diffusers.utils import load_image
-from transformers import pipeline
-
-img = load_image(
-    "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main" "/kandinskyv22/cat.png"
-).resize((768, 768))
-
-
-def make_hint(image, depth_estimator):
-    image = depth_estimator(image)["depth"]
-    image = np.array(image)
-    image = image[:, :, None]
-    image = np.concatenate([image, image, image], axis=2)
-    detected_map = torch.from_numpy(image).float() / 255.0
-    hint = detected_map.permute(2, 0, 1)
-    return hint
-
-
-depth_estimator = pipeline("depth-estimation")
-hint = make_hint(img, depth_estimator).unsqueeze(0).half().to("cuda")
-
-pipe_prior = KandinskyV22PriorEmb2EmbPipeline.from_pretrained(
-    "kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16
-)
-pipe_prior = pipe_prior.to("cuda")
-
-pipe = KandinskyV22ControlnetImg2ImgPipeline.from_pretrained(
-    "kandinsky-community/kandinsky-2-2-controlnet-depth", torch_dtype=torch.float16
-)
-pipe = pipe.to("cuda")
-
-prompt = "A robot, 4k photo"
-negative_prior_prompt = "lowres, text, error, cropped, worst quality, low quality, jpeg artifacts, ugly, duplicate, morbid, mutilated, out of frame, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, mutation, deformed, blurry, dehydrated, bad anatomy, bad proportions, extra limbs, cloned face, disfigured, gross proportions, malformed limbs, missing arms, missing legs, extra arms, extra legs, fused fingers, too many fingers, long neck, username, watermark, signature"
-
-generator = torch.Generator(device="cuda").manual_seed(43)
-
-# run prior pipeline
-
-img_emb = pipe_prior(prompt=prompt, image=img, strength=0.85, generator=generator)
-negative_emb = pipe_prior(prompt=negative_prior_prompt, image=img, strength=1, generator=generator)
-
-# run controlnet img2img pipeline
-images = pipe(
-    image=img,
-    strength=0.5,
-    image_embeds=img_emb.image_embeds,
-    negative_image_embeds=negative_emb.image_embeds,
-    hint=hint,
-    num_inference_steps=50,
-    generator=generator,
-    height=768,
-    width=768,
-).images
-
-images[0].save("robot_cat.png")
-```
-
-Here is the output. Compared with the output from our text-to-image controlnet example, it kept a lot more cat facial details from the original image and worked into the robot style we asked for.
-
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/robot_cat.png)
-
-## Kandinsky 2.2
-
-The Kandinsky 2.2 release includes robust new text-to-image models that support text-to-image generation, image-to-image generation, image interpolation, and text-guided image inpainting. The general workflow to perform these tasks using Kandinsky 2.2 is the same as in Kandinsky 2.1. First, you will need to use a prior pipeline to generate image embeddings based on your text prompt, and then use one of the image decoding pipelines to generate the output image. The only difference is that in Kandinsky 2.2, all of the decoding pipelines no longer accept the `prompt` input, and the image generation process is conditioned with only `image_embeds` and `negative_image_embeds`.
-
-Let's look at an example of how to perform text-to-image generation using Kandinsky 2.2.
-
-First, let's create the prior pipeline and text-to-image pipeline with Kandinsky 2.2 checkpoints.
-
-```python
-from diffusers import DiffusionPipeline
-import torch
-
-pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16)
-pipe_prior.to("cuda")
-
-t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16)
-t2i_pipe.to("cuda")
-```
-
-You can then use `pipe_prior` to generate image embeddings.
-
-```python
-prompt = "portrait of a women, blue eyes, cinematic"
-negative_prompt = "low quality, bad quality"
-
-image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
-```
-
-Now you can pass these embeddings to the text-to-image pipeline. When using Kandinsky 2.2 you don't need to pass the `prompt` (but you do with the previous version, Kandinsky 2.1).
-
-```
-image = t2i_pipe(image_embeds=image_embeds, negative_image_embeds=negative_image_embeds, height=768, width=768).images[
-    0
-]
-image.save("portrait.png")
-```
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/%20blue%20eyes.png)
-
-We used the text-to-image pipeline as an example, but the same process applies to all decoding pipelines in Kandinsky 2.2. For more information, please refer to our API section for each pipeline.
-
-
-## Optimization
-
-Running Kandinsky in inference requires running both a first prior pipeline: [`KandinskyPriorPipeline`]
-and a second image decoding pipeline which is one of [`KandinskyPipeline`], [`KandinskyImg2ImgPipeline`], or [`KandinskyInpaintPipeline`].
-
-The bulk of the computation time will always be the second image decoding pipeline, so when looking 
-into optimizing the model, one should look into the second image decoding pipeline.
-
-When running with PyTorch < 2.0, we strongly recommend making use of [`xformers`](https://github.com/facebookresearch/xformers)
-to speed-up the optimization. This can be done by simply running:
-
-```py
-from diffusers import DiffusionPipeline
-import torch
-
-t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
-t2i_pipe.enable_xformers_memory_efficient_attention()
-```
-
-When running on PyTorch >= 2.0, PyTorch's SDPA attention will automatically be used. For more information on 
-PyTorch's SDPA, feel free to have a look at [this blog post](https://pytorch.org/blog/accelerated-diffusers-pt-20/).
-
-To have explicit control , you can also manually set the pipeline to use PyTorch's 2.0 efficient attention:
-
-```py
-from diffusers.models.attention_processor import AttnAddedKVProcessor2_0
-
-t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor2_0())
-```
-
-The slowest and most memory intense attention processor is the default `AttnAddedKVProcessor` processor.
-We do **not** recommend using it except for testing purposes or cases where very high determistic behaviour is desired. 
-You can set it with:
-
-```py
-from diffusers.models.attention_processor import AttnAddedKVProcessor
-
-t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor())
-```
-
-With PyTorch >= 2.0, you can also use Kandinsky with `torch.compile` which depending 
-on your hardware can signficantly speed-up your inference time once the model is compiled.
-To use Kandinsksy with `torch.compile`, you can do:
-
-```py
-t2i_pipe.unet.to(memory_format=torch.channels_last)
-t2i_pipe.unet = torch.compile(t2i_pipe.unet, mode="reduce-overhead", fullgraph=True)
-```
-
-After compilation you should see a very fast inference time. For more information,
-feel free to have a look at [Our PyTorch 2.0 benchmark](https://huggingface.co/docs/diffusers/main/en/optimization/torch2.0).
-
-## Available Pipelines:
-
-| Pipeline | Tasks |
-|---|---|
-| [pipeline_kandinsky2_2.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2.py) | *Text-to-Image Generation* |
-| [pipeline_kandinsky.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky.py) | *Text-to-Image Generation* |
-| [pipeline_kandinsky2_2_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_inpaint.py) | *Image-Guided Image Generation* |
-| [pipeline_kandinsky_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_inpaint.py) | *Image-Guided Image Generation* |
-| [pipeline_kandinsky2_2_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_img2img.py) | *Image-Guided Image Generation* |
-| [pipeline_kandinsky_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_img2img.py) | *Image-Guided Image Generation* |
-| [pipeline_kandinsky2_2_controlnet.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_controlnet.py) | *Image-Guided Image Generation* |
-| [pipeline_kandinsky2_2_controlnet_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_controlnet_img2img.py) | *Image-Guided Image Generation* |
-
-
-### KandinskyV22Pipeline
-
-[[autodoc]] KandinskyV22Pipeline
-	- all
-	- __call__
-
-### KandinskyV22ControlnetPipeline
-
-[[autodoc]] KandinskyV22ControlnetPipeline
-	- all
-	- __call__
-
-### KandinskyV22ControlnetImg2ImgPipeline
-
-[[autodoc]] KandinskyV22ControlnetImg2ImgPipeline
-	- all
-	- __call__
-
-### KandinskyV22Img2ImgPipeline
-
-[[autodoc]] KandinskyV22Img2ImgPipeline
-	- all
-	- __call__
-
-### KandinskyV22InpaintPipeline
-
-[[autodoc]] KandinskyV22InpaintPipeline
-	- all
-	- __call__
-
-### KandinskyV22PriorPipeline
-
-[[autodoc]] ## KandinskyV22PriorPipeline
-	- all
-	- __call__
-	- interpolate
-
-### KandinskyV22PriorEmb2EmbPipeline
-
-[[autodoc]] KandinskyV22PriorEmb2EmbPipeline
-	- all
-	- __call__
-	- interpolate
-
-### KandinskyPriorPipeline
+## KandinskyPriorPipeline
 
 [[autodoc]] KandinskyPriorPipeline
 	- all
 	- __call__
 	- interpolate
 	
-### KandinskyPipeline
+## KandinskyPipeline
 
 [[autodoc]] KandinskyPipeline
 	- all
 	- __call__
 
-### KandinskyImg2ImgPipeline
+## KandinskyInpaintPipeline
+
+[[autodoc]] KandinskyInpaintPipeline
+	- all
+	- __call__
+	
+## KandinskyImg2ImgPipeline
 
 [[autodoc]] KandinskyImg2ImgPipeline
 	- all
 	- __call__
 
-### KandinskyInpaintPipeline
-
-[[autodoc]] KandinskyInpaintPipeline
-	- all
-	- __call__

docs/source/en/api/pipelines/overview.mdx

@@ -54,14 +54,10 @@ available a colab notebook to directly try them out.
 | [if](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
 | [if_img2img](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
 | [if_inpainting](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
-| [kandinsky](./kandinsky) | **Kandinsky** | Text-to-Image Generation | 
-| [kandinsky_inpaint](./kandinsky) | **Kandinsky** | Image-to-Image Generation | 
-| [kandinsky_img2img](./kandinsky) | **Kandinsksy** | Image-to-Image Generation |
 | [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation | 
 | [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image | 
 | [latent_diffusion_uncond](./latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation | 
 | [paint_by_example](./paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting | 
-| [paradigms](./paradigms) | [**Parallel Sampling of Diffusion Models**](https://arxiv.org/abs/2305.16317) | Text-to-Image Generation |
 | [pndm](./pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation | 
 | [score_sde_ve](./score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation | 
 | [score_sde_vp](./score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation | 
@@ -76,20 +72,21 @@ available a colab notebook to directly try them out.
 | [stable_diffusion_self_attention_guidance](./stable_diffusion/self_attention_guidance) | [**Self-Attention Guidance**](https://arxiv.org/abs/2210.00939) | Text-to-Image Generation |
 | [stable_diffusion_image_variation](./stable_diffusion/image_variation) | [**Stable Diffusion Image Variations**](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations) | Image-to-Image Generation |
 | [stable_diffusion_latent_upscale](./stable_diffusion/latent_upscale) | [**Stable Diffusion Latent Upscaler**](https://twitter.com/StabilityAI/status/1590531958815064065) | Text-Guided Super Resolution Image-to-Image |
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting | 
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Depth-to-Image Text-Guided Generation |
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
+| [stable_diffusion_2](./stable_diffusion_2/) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation | 
+| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting | 
+| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Depth-to-Image Text-Guided Generation |
+| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
 | [stable_diffusion_safe](./stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
 | [stable_unclip](./stable_unclip) | **Stable unCLIP** | Text-to-Image Generation |
 | [stable_unclip](./stable_unclip) | **Stable unCLIP** | Image-to-Image Text-Guided Generation |
 | [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
-| [text_to_video_sd](./api/pipelines/text_to_video) | [**Modelscope's Text-to-video-synthesis Model in Open Domain**](https://modelscope.cn/models/damo/text-to-video-synthesis/summary) | Text-to-Video Generation |
-| [unclip](./unclip) | [**Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation | 
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Image Variations Generation | 
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation | 
-| [vq_diffusion](./vq_diffusion) | [**Vector Quantized Diffusion Model for Text-to-Image Synthesis**](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation | 
-| [text_to_video_zero](./text_to_video_zero) | [**Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators**](https://arxiv.org/abs/2303.13439) | Text-to-Video Generation |
+| [text_to_video_sd](./api/pipelines/text_to_video) | [Modelscope's Text-to-video-synthesis Model in Open Domain](https://modelscope.cn/models/damo/text-to-video-synthesis/summary) | Text-to-Video Generation |
+| [unclip](./unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
+| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation | 
+| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation | 
+| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation | 
+| [vq_diffusion](./vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation | 
+| [text_to_video_zero](./text_to_video_zero) | [Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://arxiv.org/abs/2303.13439) | Text-to-Video Generation |
 
 
 **Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers. 

docs/source/en/api/pipelines/paradigms.mdx

@@ -1,83 +0,0 @@
-<!--Copyright 2023 ParaDiGMS authors and The HuggingFace Team. All rights reserved.
-
-Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
-the License. You may obtain a copy of the License at
-
-http://www.apache.org/licenses/LICENSE-2.0
-
-Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
-an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
-specific language governing permissions and limitations under the License.
--->
-
-# Parallel Sampling of Diffusion Models (ParaDiGMS)
-
-## Overview
-
-[Parallel Sampling of Diffusion Models](https://arxiv.org/abs/2305.16317) by Andy Shih, Suneel Belkhale, Stefano Ermon, Dorsa Sadigh, Nima Anari.
-
-The abstract of the paper is the following:
-
-*Diffusion models are powerful generative models but suffer from slow sampling, often taking 1000 sequential denoising steps for one sample. As a result, considerable efforts have been directed toward reducing the number of denoising steps, but these methods hurt sample quality. Instead of reducing the number of denoising steps (trading quality for speed), in this paper we explore an orthogonal approach: can we run the denoising steps in parallel (trading compute for speed)? In spite of the sequential nature of the denoising steps, we show that surprisingly it is possible to parallelize sampling via Picard iterations, by guessing the solution of future denoising steps and iteratively refining until convergence. With this insight, we present ParaDiGMS, a novel method to accelerate the sampling of pretrained diffusion models by denoising multiple steps in parallel. ParaDiGMS is the first diffusion sampling method that enables trading compute for speed and is even compatible with existing fast sampling techniques such as DDIM and DPMSolver. Using ParaDiGMS, we improve sampling speed by 2-4x across a range of robotics and image generation models, giving state-of-the-art sampling speeds of 0.2s on 100-step DiffusionPolicy and 16s on 1000-step StableDiffusion-v2 with no measurable degradation of task reward, FID score, or CLIP score.*
-
-Resources:
-
-* [Paper](https://arxiv.org/abs/2305.16317).
-* [Original Code](https://github.com/AndyShih12/paradigms).
-
-## Available Pipelines:
-
-| Pipeline | Tasks | Demo
-|---|---|:---:|
-| [StableDiffusionParadigmsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_paradigms.py) | *Faster Text-to-Image Generation* | |
-
-This pipeline was contributed by [`AndyShih12`](https://github.com/AndyShih12) in this [PR](https://github.com/huggingface/diffusers/pull/3716/).
-
-## Usage example
-
-```python
-import torch
-from diffusers import DDPMParallelScheduler
-from diffusers import StableDiffusionParadigmsPipeline
-
-scheduler = DDPMParallelScheduler.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
-
-pipe = StableDiffusionParadigmsPipeline.from_pretrained(
-    "runwayml/stable-diffusion-v1-5", scheduler=scheduler, torch_dtype=torch.float16
-)
-pipe = pipe.to("cuda")
-
-ngpu, batch_per_device = torch.cuda.device_count(), 5
-pipe.wrapped_unet = torch.nn.DataParallel(pipe.unet, device_ids=[d for d in range(ngpu)])
-
-prompt = "a photo of an astronaut riding a horse on mars"
-image = pipe(prompt, parallel=ngpu * batch_per_device, num_inference_steps=1000).images[0]
-```
-
-<Tip>
-This pipeline improves sampling speed by running denoising steps in parallel, at the cost of increased total FLOPs.
-Therefore, it is better to call this pipeline when running on multiple GPUs. Otherwise, without enough GPU bandwidth
-sampling may be even slower than sequential sampling.
-
-The two parameters to play with are `parallel` (batch size) and `tolerance`. 
-- If it fits in memory, for 1000-step DDPM you can aim for a batch size of around 100 
-(e.g. 8 GPUs and batch_per_device=12 to get parallel=96). Higher batch size
-may not fit in memory, and lower batch size gives less parallelism. 
-- For tolerance, using a higher tolerance may get better speedups but can risk sample quality degradation. 
-If there is quality degradation with the default tolerance, then use a lower tolerance (e.g. 0.001).
-
-For 1000-step DDPM on 8 A100 GPUs, you can expect around a 3x speedup by StableDiffusionParadigmsPipeline instead of StableDiffusionPipeline
-by setting parallel=80 and tolerance=0.1.
-</Tip>
-
-<Tip>
-Diffusers also offers distributed inference support for generating multiple prompts
-in parallel on multiple GPUs. Check out the docs [here](https://huggingface.co/docs/diffusers/main/en/training/distributed_inference).
-
-In contrast, this pipeline is designed for speeding up sampling of a single prompt (by using multiple GPUs).
-</Tip>
-
-## StableDiffusionParadigmsPipeline
-[[autodoc]] StableDiffusionParadigmsPipeline
-	- __call__
-	- all

docs/source/en/api/pipelines/shap_e.mdx

@@ -1,139 +0,0 @@
-<!--Copyright 2023 The HuggingFace Team. All rights reserved.
-Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
-the License. You may obtain a copy of the License at
-http://www.apache.org/licenses/LICENSE-2.0
-Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
-an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
-specific language governing permissions and limitations under the License.
--->
-
-# Shap-E
-
-## Overview
-
-
-The Shap-E model was proposed in [Shap-E: Generating Conditional 3D Implicit Functions](https://arxiv.org/abs/2305.02463) by Alex Nichol and Heewon Jun from [OpenAI](https://github.com/openai). 
-
-The abstract of the paper is the following:
-
-*We present Shap-E, a conditional generative model for 3D assets. Unlike recent work on 3D generative models which produce a single output representation, Shap-E directly generates the parameters of implicit functions that can be rendered as both textured meshes and neural radiance fields. We train Shap-E in two stages: first, we train an encoder that deterministically maps 3D assets into the parameters of an implicit function; second, we train a conditional diffusion model on outputs of the encoder. When trained on a large dataset of paired 3D and text data, our resulting models are capable of generating complex and diverse 3D assets in a matter of seconds. When compared to Point-E, an explicit generative model over point clouds, Shap-E converges faster and reaches comparable or better sample quality despite modeling a higher-dimensional, multi-representation output space.*
-
-The original codebase can be found [here](https://github.com/openai/shap-e).
-
-## Available Pipelines:
-
-| Pipeline | Tasks |
-|---|---|
-| [pipeline_shap_e.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/shap_e/pipeline_shap_e.py) | *Text-to-Image Generation* | 
-| [pipeline_shap_e_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/shap_e/pipeline_shap_e_img2img.py) | *Image-to-Image Generation* |
-
-## Available checkpoints 
-
-* [`openai/shap-e`](https://huggingface.co/openai/shap-e)
-* [`openai/shap-e-img2img`](https://huggingface.co/openai/shap-e-img2img)
-
-## Usage Examples
-
-In the following, we will walk you through some examples of how to use Shap-E pipelines to create 3D objects in gif format.
-
-### Text-to-3D image generation 
-
-We can use [`ShapEPipeline`] to create 3D object based on a text prompt. In this example, we will make a birthday cupcake for :firecracker: diffusers library's 1 year birthday. The workflow to use the Shap-E text-to-image pipeline is same as how you would use other text-to-image pipelines in diffusers.
-
-```python
-import torch
-
-from diffusers import DiffusionPipeline
-
-device = torch.device("cuda" if torch.cuda.is_available() else "cpu")
-
-repo = "openai/shap-e"
-pipe = DiffusionPipeline.from_pretrained(repo, torch_dtype=torch.float16)
-pipe = pipe.to(device)
-
-guidance_scale = 15.0
-prompt = ["A firecracker", "A birthday cupcake"]
-
-images = pipe(
-    prompt,
-    guidance_scale=guidance_scale,
-    num_inference_steps=64,
-    frame_size=256,
-).images
-```
-
-The output of [`ShapEPipeline`] is a list of lists of images frames. Each list of frames can be used to create a 3D object. Let's use the `export_to_gif` utility function in diffusers to make a 3D cupcake!
-
-```python
-from diffusers.utils import export_to_gif
-
-export_to_gif(images[0], "firecracker_3d.gif")
-export_to_gif(images[1], "cake_3d.gif")
-```
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/firecracker_out.gif)
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/cake_out.gif)
-
-
-### Image-to-Image generation
-
-You can use [`ShapEImg2ImgPipeline`] along with other text-to-image pipelines in diffusers and turn your 2D generation into 3D. 
-
-In this example, We will first genrate a cheeseburger with a simple prompt "A cheeseburger, white background" 
-
-```python
-from diffusers import DiffusionPipeline
-import torch
-
-pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16)
-pipe_prior.to("cuda")
-
-t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
-t2i_pipe.to("cuda")
-
-prompt = "A cheeseburger, white background"
-
-image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
-image = t2i_pipe(
-    prompt,
-    image_embeds=image_embeds,
-    negative_image_embeds=negative_image_embeds,
-).images[0]
-
-image.save("burger.png")
-```
-
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/burger_in.png)
-
-we will then use the Shap-E image-to-image pipeline to turn it into a 3D cheeseburger :)
-
-```python
-from PIL import Image
-from diffusers.utils import export_to_gif
-
-repo = "openai/shap-e-img2img"
-pipe = DiffusionPipeline.from_pretrained(repo, torch_dtype=torch.float16)
-pipe = pipe.to("cuda")
-
-guidance_scale = 3.0
-image = Image.open("burger.png").resize((256, 256))
-
-images = pipe(
-    image,
-    guidance_scale=guidance_scale,
-    num_inference_steps=64,
-    frame_size=256,
-).images
-
-gif_path = export_to_gif(images[0], "burger_3d.gif")
-```
-![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/burger_out.gif)
-
-## ShapEPipeline
-[[autodoc]] ShapEPipeline
-	- all
-	- __call__
-
-## ShapEImg2ImgPipeline
-[[autodoc]] ShapEImg2ImgPipeline
-	- all
-	- __call__

docs/source/en/api/pipelines/stable_diffusion/img2img.mdx

@@ -31,7 +31,7 @@ proposed by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan
 	- enable_xformers_memory_efficient_attention
 	- disable_xformers_memory_efficient_attention
 	- load_textual_inversion
-	- from_single_file
+	- from_ckpt
 	- load_lora_weights
 	- save_lora_weights
 

docs/source/en/api/pipelines/stable_diffusion/ldm3d_diffusion.mdx

@@ -1,55 +0,0 @@
-<!--Copyright 2023 The Intel Labs Team Authors and HuggingFace Team. All rights reserved.
-
-Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
-the License. You may obtain a copy of the License at
-
-http://www.apache.org/licenses/LICENSE-2.0
-
-Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
-an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
-specific language governing permissions and limitations under the License.
--->
-
-# LDM3D
-
-LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://arxiv.org/abs/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, Vasudev Lal
-The abstract of the paper is the following:
-
-*This research paper proposes a Latent Diffusion Model for 3D (LDM3D) that generates both image and depth map data from a given text prompt, allowing users to generate RGBD images from text prompts. The LDM3D model is fine-tuned on a dataset of tuples containing an RGB image, depth map and caption, and validated through extensive experiments. We also develop an application called DepthFusion, which uses the generated RGB images and depth maps to create immersive and interactive 360-degree-view experiences using TouchDesigner. This technology has the potential to transform a wide range of industries, from entertainment and gaming to architecture and design. Overall, this paper presents a significant contribution to the field of generative AI and computer vision, and showcases the potential of LDM3D and DepthFusion to revolutionize content creation and digital experiences. A short video summarizing the approach can be found at [this url](https://t.ly/tdi2).*
-
-
-*Overview*:
-
-| Pipeline | Tasks | Colab | Demo
-|---|---|:---:|:---:|
-| [pipeline_stable_diffusion_ldm3d.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_ldm3d.py) | *Text-to-Image Generation* | - | -
-
-## Tips
-
-- LDM3D generates both an image and a depth map from a given text prompt, compared to the existing txt-to-img diffusion models such as [Stable Diffusion](./stable_diffusion/overview) that generates only an image.
-- With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps. 
-
-
-Running LDM3D is straighforward with the [`StableDiffusionLDM3DPipeline`]:
-
-```python
->>> from diffusers import StableDiffusionLDM3DPipeline
-
->>> pipe = StableDiffusionLDM3DPipeline.from_pretrained("Intel/ldm3d")
-prompt ="A picture of some lemons on a table"
-output = pipe(prompt)
-rgb_image, depth_image = output.rgb, output.depth
-rgb_image[0].save("lemons_ldm3d_rgb.jpg")
-depth_image[0].save("lemons_ldm3d_depth.png")
-```
-
-
-## StableDiffusionPipelineOutput
-[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
-	- all
-	- __call__
-
-## StableDiffusionLDM3DPipeline
-[[autodoc]] StableDiffusionLDM3DPipeline
-	- all
-	- __call__

docs/source/en/api/pipelines/stable_diffusion/overview.mdx

@@ -26,17 +26,19 @@ For more details about how Stable Diffusion works and how it differs from the ba
 | Pipeline | Tasks | Colab | Demo
 |---|---|:---:|:---:|
 | [StableDiffusionPipeline](./text2img) | *Text-to-Image Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) | [🤗 Stable Diffusion](https://huggingface.co/spaces/stabilityai/stable-diffusion)
-| [StableDiffusionPipelineSafe](./stable_diffusion_safe) | *Text-to-Image Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb) | [![Huggingface Spaces](https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue)](https://huggingface.co/spaces/AIML-TUDA/unsafe-vs-safe-stable-diffusion)
 | [StableDiffusionImg2ImgPipeline](./img2img) | *Image-to-Image Text-Guided Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) | [🤗 Diffuse the Rest](https://huggingface.co/spaces/huggingface/diffuse-the-rest)
-| [StableDiffusionInpaintPipeline](./inpaint) | **Experimental** – *Text-Guided Image Inpainting* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | 
-| [StableDiffusionDepth2ImgPipeline](./depth2img) | **Experimental** – *Depth-to-Image Text-Guided Generation* | |
-| [StableDiffusionImageVariationPipeline](./image_variation) | **Experimental** – *Image Variation Generation* | | [🤗 Stable Diffusion Image Variations](https://huggingface.co/spaces/lambdalabs/stable-diffusion-image-variations)
-| [StableDiffusionUpscalePipeline](./upscale) | **Experimental** – *Text-Guided Image Super-Resolution* | |
-| [StableDiffusionLatentUpscalePipeline](./latent_upscale) | **Experimental** – *Text-Guided Image Super-Resolution* | |
-| [Stable Diffusion 2](./stable_diffusion_2)  | *Text-Guided Image Inpainting* | 
-| [Stable Diffusion 2](./stable_diffusion_2)  |  *Depth-to-Image Text-Guided Generation* |
-| [Stable Diffusion 2](./stable_diffusion_2)  | *Text-Guided Super Resolution Image-to-Image* |
-| [StableDiffusionLDM3DPipeline](./ldm3d)  | *Text-to-(RGB, Depth)* |
+| [StableDiffusionInpaintPipeline](./inpaint) | **Experimental** – *Text-Guided Image Inpainting* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | Coming soon
+| [StableDiffusionDepth2ImgPipeline](./depth2img) | **Experimental** – *Depth-to-Image Text-Guided Generation * | | Coming soon
+| [StableDiffusionImageVariationPipeline](./image_variation) | **Experimental** – *Image Variation Generation * | | [🤗 Stable Diffusion Image Variations](https://huggingface.co/spaces/lambdalabs/stable-diffusion-image-variations)
+| [StableDiffusionUpscalePipeline](./upscale) | **Experimental** – *Text-Guided Image Super-Resolution * | | Coming soon
+| [StableDiffusionLatentUpscalePipeline](./latent_upscale) | **Experimental** – *Text-Guided Image Super-Resolution * | | Coming soon
+| [StableDiffusionInstructPix2PixPipeline](./pix2pix) | **Experimental** – *Text-Based Image Editing * | | [InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/spaces/timbrooks/instruct-pix2pix)
+| [StableDiffusionAttendAndExcitePipeline](./attend_and_excite) | **Experimental** – *Text-to-Image Generation * | | [Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite)
+| [StableDiffusionPix2PixZeroPipeline](./pix2pix_zero) | **Experimental** – *Text-Based Image Editing * | | [Zero-shot Image-to-Image Translation](https://arxiv.org/abs/2302.03027)
+| [StableDiffusionModelEditingPipeline](./model_editing) | **Experimental** – *Text-to-Image Model Editing * | | [Editing Implicit Assumptions in Text-to-Image Diffusion Models](https://arxiv.org/abs/2303.08084)
+| [StableDiffusionDiffEditPipeline](./diffedit) | **Experimental** – *Text-Based Image Editing * | | [DiffEdit: Diffusion-based semantic image editing with mask guidance](https://arxiv.org/abs/2210.11427)
+
+
 
 ## Tips
 

docs/source/en/api/pipelines/stable_diffusion/panorama.mdx

@@ -52,14 +52,6 @@ image = pipe(prompt).images[0]
 image.save("dolomites.png")
 ```
 
-<Tip>
-
-While calling this pipeline, it's possible to specify the `view_batch_size` to have a >1 value. 
-For some GPUs with high performance, higher a `view_batch_size`, can speedup the generation
-and increase the VRAM usage.
-
-</Tip>
-
 ## StableDiffusionPanoramaPipeline
 [[autodoc]] StableDiffusionPanoramaPipeline
 	- __call__

docs/source/en/api/pipelines/stable_diffusion/stable_diffusion_xl.mdx

@@ -1,177 +0,0 @@
-<!--Copyright 2023 The HuggingFace Team. All rights reserved.
-
-Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
-the License. You may obtain a copy of the License at
-
-http://www.apache.org/licenses/LICENSE-2.0
-
-Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
-an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
-specific language governing permissions and limitations under the License.
--->
-
-# Stable diffusion XL
-
-Stable Diffusion XL was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://arxiv.org/abs/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, Robin Rombach
-
-The abstract of the paper is the following:
-
-*We present SDXL, a latent diffusion model for text-to-image synthesis. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. We design multiple novel conditioning schemes and train SDXL on multiple aspect ratios. We also introduce a refinement model which is used to improve the visual fidelity of samples generated by SDXL using a post-hoc image-to-image technique. We demonstrate that SDXL shows drastically improved performance compared the previous versions of Stable Diffusion and achieves results competitive with those of black-box state-of-the-art image generators.*
-
-## Tips
-
-- Stable Diffusion XL works especially well with images between 768 and 1024.
-- Stable Diffusion XL output image can be improved by making use of a refiner as shown below.
-
-### Available checkpoints:
-
-- *Text-to-Image (1024x1024 resolution)*: [stabilityai/stable-diffusion-xl-base-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9) with [`StableDiffusionXLPipeline`]
-- *Image-to-Image / Refiner (1024x1024 resolution)*: [stabilityai/stable-diffusion-xl-refiner-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9) with [`StableDiffusionXLImg2ImgPipeline`]
-
-## Usage Example
-
-Before using SDXL make sure to have `transformers`, `accelerate`, `safetensors` and `invisible_watermark` installed. 
-You can install the libraries as follows:
-
-```
-pip install transformers
-pip install accelerate
-pip install safetensors
-pip install invisible-watermark>=2.0
-```
-
-### Text-to-Image
-
-You can use SDXL as follows for *text-to-image*:
-
-```py
-from diffusers import StableDiffusionXLPipeline
-import torch
-
-pipe = StableDiffusionXLPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-0.9", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
-)
-pipe.to("cuda")
-
-prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
-image = pipe(prompt=prompt).images[0]
-```
-
-### Refining the image output
-
-The image can be refined by making use of [stabilityai/stable-diffusion-xl-refiner-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9).
-In this case, you only have to output the `latents` from the base model.
-
-```py
-from diffusers import StableDiffusionXLPipeline, StableDiffusionXLImg2ImgPipeline
-import torch
-
-pipe = StableDiffusionXLPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-0.9", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
-)
-pipe.to("cuda")
-
-use_refiner = True
-refiner = StableDiffusionXLImg2ImgPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-refiner-0.9", torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
-)
-refiner.to("cuda")
-
-prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
-
-image = pipe(prompt=prompt, output_type="latent" if use_refiner else "pil").images[0]
-image = refiner(prompt=prompt, image=image[None, :]).images[0]
-```
-
-### Image-to-image 
-
-```py 
-import torch
-from diffusers import StableDiffusionXLImg2ImgPipeline
-from diffusers.utils import load_image
-
-pipe = StableDiffusionXLImg2ImgPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-refiner-0.9", torch_dtype=torch.float16
-)
-pipe = pipe.to("cuda")
-url = "https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/aa_xl/000000009.png"
-
-init_image = load_image(url).convert("RGB")
-prompt = "a photo of an astronaut riding a horse on mars"
-image = pipe(prompt, image=init_image).images[0]
-```
-
-| Original Image | Refined Image |
-|---|---|
-| ![](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/sd_xl/init_image.png) | ![](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/sd_xl/refined_image.png) |
-
-### Loading single file checkpoints / original file format
-
-By making use of [`~diffusers.loaders.FromSingleFileMixin.from_single_file`] you can also load the 
-original file format into `diffusers`:
-
-```py
-from diffusers import StableDiffusionXLPipeline, StableDiffusionXLImg2ImgPipeline
-import torch
-
-pipe = StableDiffusionXLPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-0.9", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
-)
-pipe.to("cuda")
-
-refiner = StableDiffusionXLImg2ImgPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-refiner-0.9", torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
-)
-refiner.to("cuda")
-```
-
-### Memory optimization via model offloading 
-
-If you are seeing out-of-memory errors, we recommend making use of [`StableDiffusionXLPipeline.enable_model_cpu_offload`].
-
-```diff
-- pipe.to("cuda")
-+ pipe.enable_model_cpu_offload()
-```
-
-and 
-
-```diff
-- refiner.to("cuda")
-+ refiner.enable_model_cpu_offload()
-```
-
-### Speed-up inference with `torch.compile`
-
-You can speed up inference by making use of `torch.compile`. This should give you **ca.** 20% speed-up.
-
-```diff
-+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
-+ refiner.unet = torch.compile(refiner.unet, mode="reduce-overhead", fullgraph=True)
-```
-
-### Running with `torch` \< 2.0
-
-**Note** that if you want to run Stable Diffusion XL with `torch` < 2.0, please make sure to enable xformers 
-attention:
-
-```
-pip install xformers
-```
-
-```diff
-+pipe.enable_xformers_memory_efficient_attention()
-+refiner.enable_xformers_memory_efficient_attention()
-```
-
-## StableDiffusionXLPipeline
-
-[[autodoc]] StableDiffusionXLPipeline
-	- all
-	- __call__
-
-## StableDiffusionXLImg2ImgPipeline
-
-[[autodoc]] StableDiffusionXLImg2ImgPipeline
-	- all
-	- __call__

docs/source/en/api/pipelines/stable_diffusion/text2img.mdx

@@ -40,7 +40,7 @@ Available Checkpoints are:
 	- enable_vae_tiling
 	- disable_vae_tiling
 	- load_textual_inversion
-	- from_single_file
+	- from_ckpt
 	- load_lora_weights
 	- save_lora_weights
 

docs/source/en/api/pipelines/stable_diffusion_2.mdx

@@ -71,64 +71,6 @@ image = pipe(prompt, guidance_scale=9, num_inference_steps=25).images[0]
 image.save("astronaut.png")
 ```
 
-#### Experimental: "Common Diffusion Noise Schedules and Sample Steps are Flawed": 
-
-The paper **[Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/abs/2305.08891)** 
-claims that a mismatch between the training and inference settings leads to suboptimal inference generation results for Stable Diffusion.
-
-The abstract reads as follows:
-
-*We discover that common diffusion noise schedules do not enforce the last timestep to have zero signal-to-noise ratio (SNR),
-and some implementations of diffusion samplers do not start from the last timestep.
-Such designs are flawed and do not reflect the fact that the model is given pure Gaussian noise at inference, creating a discrepancy between training and inference.
-We show that the flawed design causes real problems in existing implementations. 
-In Stable Diffusion, it severely limits the model to only generate images with medium brightness and 
-prevents it from generating very bright and dark samples. We propose a few simple fixes: 
-- (1) rescale the noise schedule to enforce zero terminal SNR; 
-- (2) train the model with v prediction; 
-- (3) change the sampler to always start from the last timestep; 
-- (4) rescale classifier-free guidance to prevent over-exposure. 
-These simple changes ensure the diffusion process is congruent between training and inference and 
-allow the model to generate samples more faithful to the original data distribution.*
-
-You can apply all of these changes in `diffusers` when using [`DDIMScheduler`]:
-- (1) rescale the noise schedule to enforce zero terminal SNR; 
-```py
-pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, rescale_betas_zero_snr=True)
-```
-- (2) train the model with v prediction; 
-Continue fine-tuning a checkpoint with [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py)
-and `--prediction_type="v_prediction"`.
-- (3) change the sampler to always start from the last timestep; 
-```py
-pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, timestep_spacing="trailing")
-```
-- (4) rescale classifier-free guidance to prevent over-exposure. 
-```py
-pipe(..., guidance_rescale=0.7)
-```
-
-An example is to use [this checkpoint](https://huggingface.co/ptx0/pseudo-journey-v2) 
-which has been fine-tuned using the `"v_prediction"`.
-
-The checkpoint can then be run in inference as follows:
-
-```py
-from diffusers import DiffusionPipeline, DDIMScheduler
-
-pipe = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2", torch_dtype=torch.float16)
-pipe.scheduler = DDIMScheduler.from_config(
-    pipe.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
-)
-pipe.to("cuda")
-
-prompt = "A lion in galaxies, spirals, nebulae, stars, smoke, iridescent, intricate detail, octane render, 8k"
-image = pipeline(prompt, guidance_rescale=0.7).images[0]
-```
-
-## DDIMScheduler
-[[autodoc]] DDIMScheduler
-
 ### Image Inpainting
 
 - *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]

docs/source/en/api/pipelines/text_to_video.mdx

@@ -37,12 +37,9 @@ Resources:
 | Pipeline | Tasks | Demo
 |---|---|:---:|
 | [TextToVideoSDPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/text_to_video_synthesis/pipeline_text_to_video_synth.py) | *Text-to-Video Generation* | [🤗 Spaces](https://huggingface.co/spaces/damo-vilab/modelscope-text-to-video-synthesis)
-| [VideoToVideoSDPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/text_to_video_synthesis/pipeline_text_to_video_synth_img2img.py) | *Text-Guided Video-to-Video Generation* | [(TODO)🤗 Spaces]()
 
 ## Usage example 
 
-### `text-to-video-ms-1.7b`
-
 Let's start by generating a short video with the default length of 16 frames (2s at 8 fps):
 
 ```python 
@@ -122,98 +119,12 @@ Here are some sample outputs:
     </tr>
 </table>
 
-### `cerspense/zeroscope_v2_576w` & `cerspense/zeroscope_v2_XL`
-
-Zeroscope are watermark-free model and have been trained on specific sizes such as `576x320` and `1024x576`.
-One should first generate a video using the lower resolution checkpoint [`cerspense/zeroscope_v2_576w`](https://huggingface.co/cerspense/zeroscope_v2_576w) with [`TextToVideoSDPipeline`],
-which can then be upscaled using [`VideoToVideoSDPipeline`] and [`cerspense/zeroscope_v2_XL`](https://huggingface.co/cerspense/zeroscope_v2_XL).
-
-
-```py
-import torch
-from diffusers import DiffusionPipeline
-from diffusers.utils import export_to_video
-
-pipe = DiffusionPipeline.from_pretrained("cerspense/zeroscope_v2_576w", torch_dtype=torch.float16)
-pipe.enable_model_cpu_offload()
-
-# memory optimization
-pipe.unet.enable_forward_chunking(chunk_size=1, dim=1)
-pipe.enable_vae_slicing()
-
-prompt = "Darth Vader surfing a wave"
-video_frames = pipe(prompt, num_frames=24).frames
-video_path = export_to_video(video_frames)
-video_path
-```
-
-Now the video can be upscaled:
-
-```py
-pipe = DiffusionPipeline.from_pretrained("cerspense/zeroscope_v2_XL", torch_dtype=torch.float16)
-pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
-pipe.enable_model_cpu_offload()
-
-# memory optimization
-pipe.unet.enable_forward_chunking(chunk_size=1, dim=1)
-pipe.enable_vae_slicing()
-
-video = [Image.fromarray(frame).resize((1024, 576)) for frame in video_frames]
-
-video_frames = pipe(prompt, video=video, strength=0.6).frames
-video_path = export_to_video(video_frames)
-video_path
-```
-
-Here are some sample outputs: 
-
-<table>
-    <tr>
-        <td ><center>
-        Darth vader surfing in waves.
-        <br>
-        <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/darthvader_cerpense.gif"
-            alt="Darth vader surfing in waves."
-            style="width: 576px;" />
-        </center></td>
-    </tr>
-</table>
-
-### Memory optimizations
-
-Text-guided video generation with [`~TextToVideoSDPipeline`] and [`~VideoToVideoSDPipeline`] is very memory intensive both
-when denoising with [`~UNet3DConditionModel`] and when decoding with [`~AutoencoderKL`]. It is possible though to reduce 
-memory usage at the cost of increased runtime to achieve the exact same result. To do so, it is recommended to enable
-**forward chunking** and **vae slicing**:
-
-Forward chunking via [`~UNet3DConditionModel.enable_forward_chunking`]is explained in [this blog post](https://huggingface.co/blog/reformer#2-chunked-feed-forward-layers) and 
-allows to significantly reduce the required memory for the unet. You can chunk the feed forward layer over the `num_frames`
-dimension by doing:
-
-```py
-pipe.unet.enable_forward_chunking(chunk_size=1, dim=1)
-```
-
-Vae slicing via [`~TextToVideoSDPipeline.enable_vae_slicing`] and [`~VideoToVideoSDPipeline.enable_vae_slicing`] also 
-gives significant memory savings since the two pipelines decode all image frames at once.
-
-```py
-pipe.enable_vae_slicing()
-```
-
 ## Available checkpoints 
 
 * [damo-vilab/text-to-video-ms-1.7b](https://huggingface.co/damo-vilab/text-to-video-ms-1.7b/)
 * [damo-vilab/text-to-video-ms-1.7b-legacy](https://huggingface.co/damo-vilab/text-to-video-ms-1.7b-legacy)
-* [cerspense/zeroscope_v2_576w](https://huggingface.co/cerspense/zeroscope_v2_576w)
-* [cerspense/zeroscope_v2_XL](https://huggingface.co/cerspense/zeroscope_v2_XL)
 
 ## TextToVideoSDPipeline
 [[autodoc]] TextToVideoSDPipeline
 	- all
 	- __call__
-
-## VideoToVideoSDPipeline
-[[autodoc]] VideoToVideoSDPipeline
-	- all
-	- __call__

docs/source/en/api/pipelines/text_to_video_zero.mdx

@@ -80,41 +80,6 @@ You can change these parameters in the pipeline call:
 * Video length:
     * `video_length`, the number of frames video_length to be generated. Default: `video_length=8`
 
-We an also generate longer videos by doing the processing in a chunk-by-chunk manner:
-```python
-import torch
-import imageio
-from diffusers import TextToVideoZeroPipeline
-import numpy as np
-
-model_id = "runwayml/stable-diffusion-v1-5"
-pipe = TextToVideoZeroPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
-seed = 0
-video_length = 8
-chunk_size = 4
-prompt = "A panda is playing guitar on times square"
-
-# Generate the video chunk-by-chunk
-result = []
-chunk_ids = np.arange(0, video_length, chunk_size - 1)
-generator = torch.Generator(device="cuda")
-for i in range(len(chunk_ids)):
-    print(f"Processing chunk {i + 1} / {len(chunk_ids)}")
-    ch_start = chunk_ids[i]
-    ch_end = video_length if i == len(chunk_ids) - 1 else chunk_ids[i + 1]
-    # Attach the first frame for Cross Frame Attention
-    frame_ids = [0] + list(range(ch_start, ch_end))
-    # Fix the seed for the temporal consistency
-    generator.manual_seed(seed)
-    output = pipe(prompt=prompt, video_length=len(frame_ids), generator=generator, frame_ids=frame_ids)
-    result.append(output.images[1:])
-
-# Concatenate chunks and save
-result = np.concatenate(result)
-result = [(r * 255).astype("uint8") for r in result]
-imageio.mimsave("video.mp4", result, fps=4)
-```
-
 
 ### Text-To-Video with Pose Control
 To generate a video from prompt with additional pose control
@@ -237,7 +202,7 @@ can run with custom [DreamBooth](../training/dreambooth) models, as shown below
 
     reader = imageio.get_reader(video_path, "ffmpeg")
     frame_count = 8
-    canny_edges = [Image.fromarray(reader.get_data(i)) for i in range(frame_count)]
+    video = [Image.fromarray(reader.get_data(i)) for i in range(frame_count)]
     ```
 
 3. Run `StableDiffusionControlNetPipeline` with custom trained DreamBooth model
@@ -258,10 +223,10 @@ can run with custom [DreamBooth](../training/dreambooth) models, as shown below
     pipe.controlnet.set_attn_processor(CrossFrameAttnProcessor(batch_size=2))
 
     # fix latents for all frames
-    latents = torch.randn((1, 4, 64, 64), device="cuda", dtype=torch.float16).repeat(len(canny_edges), 1, 1, 1)
+    latents = torch.randn((1, 4, 64, 64), device="cuda", dtype=torch.float16).repeat(len(pose_images), 1, 1, 1)
 
     prompt = "oil painting of a beautiful girl avatar style"
-    result = pipe(prompt=[prompt] * len(canny_edges), image=canny_edges, latents=latents).images
+    result = pipe(prompt=[prompt] * len(pose_images), image=pose_images, latents=latents).images
     imageio.mimsave("video.mp4", result, fps=4)
     ```
 

docs/source/en/api/schedulers/cm_stochastic_iterative.mdx

@@ -1,11 +0,0 @@
-# Consistency Model Multistep Scheduler
-
-## Overview
-
-Multistep and onestep scheduler (Algorithm 1) introduced alongside consistency models in the paper [Consistency Models](https://arxiv.org/abs/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
-Based on the [original consistency models implementation](https://github.com/openai/consistency_models).
-Should generate good samples from [`ConsistencyModelPipeline`] in one or a small number of steps.
-
-## CMStochasticIterativeScheduler
-[[autodoc]] CMStochasticIterativeScheduler
-

docs/source/en/api/schedulers/ddim.mdx

@@ -18,71 +18,10 @@ specific language governing permissions and limitations under the License.
 
 The abstract of the paper is the following:
 
-*Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, 
-yet they require simulating a Markov chain for many steps to produce a sample. 
-To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models
-with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. 
-We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from.
-We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off 
-computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.*
+Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
 
 The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
 For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
 
-### Experimental: "Common Diffusion Noise Schedules and Sample Steps are Flawed": 
-
-The paper **[Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/abs/2305.08891)** 
-claims that a mismatch between the training and inference settings leads to suboptimal inference generation results for Stable Diffusion.
-
-The abstract reads as follows:
-
-*We discover that common diffusion noise schedules do not enforce the last timestep to have zero signal-to-noise ratio (SNR),
-and some implementations of diffusion samplers do not start from the last timestep.
-Such designs are flawed and do not reflect the fact that the model is given pure Gaussian noise at inference, creating a discrepancy between training and inference.
-We show that the flawed design causes real problems in existing implementations. 
-In Stable Diffusion, it severely limits the model to only generate images with medium brightness and 
-prevents it from generating very bright and dark samples. We propose a few simple fixes: 
-- (1) rescale the noise schedule to enforce zero terminal SNR; 
-- (2) train the model with v prediction; 
-- (3) change the sampler to always start from the last timestep; 
-- (4) rescale classifier-free guidance to prevent over-exposure. 
-These simple changes ensure the diffusion process is congruent between training and inference and 
-allow the model to generate samples more faithful to the original data distribution.*
-
-You can apply all of these changes in `diffusers` when using [`DDIMScheduler`]:
-- (1) rescale the noise schedule to enforce zero terminal SNR; 
-```py
-pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, rescale_betas_zero_snr=True)
-```
-- (2) train the model with v prediction; 
-Continue fine-tuning a checkpoint with [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py)
-and `--prediction_type="v_prediction"`.
-- (3) change the sampler to always start from the last timestep; 
-```py
-pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, timestep_spacing="trailing")
-```
-- (4) rescale classifier-free guidance to prevent over-exposure. 
-```py
-pipe(..., guidance_rescale=0.7)
-```
-
-An example is to use [this checkpoint](https://huggingface.co/ptx0/pseudo-journey-v2) 
-which has been fine-tuned using the `"v_prediction"`.
-
-The checkpoint can then be run in inference as follows:
-
-```py
-from diffusers import DiffusionPipeline, DDIMScheduler
-
-pipe = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2", torch_dtype=torch.float16)
-pipe.scheduler = DDIMScheduler.from_config(
-    pipe.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
-)
-pipe.to("cuda")
-
-prompt = "A lion in galaxies, spirals, nebulae, stars, smoke, iridescent, intricate detail, octane render, 8k"
-image = pipeline(prompt, guidance_rescale=0.7).images[0]
-```
-
 ## DDIMScheduler
 [[autodoc]] DDIMScheduler

docs/source/en/api/utilities.mdx

@@ -1,23 +0,0 @@
-# Utilities
-
-Utility and helper functions for working with 🤗 Diffusers.
-
-## randn_tensor
-
-[[autodoc]] diffusers.utils.randn_tensor
-
-## numpy_to_pil
-
-[[autodoc]] utils.pil_utils.numpy_to_pil
-
-## pt_to_pil
-
-[[autodoc]] utils.pil_utils.pt_to_pil
-
-## load_image
-
-[[autodoc]] utils.testing_utils.load_image
-
-## export_to_video
-
-[[autodoc]] utils.testing_utils.export_to_video

docs/source/en/index.mdx

@@ -94,4 +94,3 @@ The library has three main components:
 | [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
 | [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
 | [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
-| [stable_diffusion_ldm3d](./api/pipelines/stable_diffusion/ldm3d_diffusion) | [LDM3D: Latent Diffusion Model for 3D](https://arxiv.org/abs/2305.10853) | Text to Image and Depth Generation |

docs/source/en/installation.mdx

@@ -23,7 +23,7 @@ Install 🤗 Diffusers for whichever deep learning library you're working with.
 
 You should install 🤗 Diffusers in a [virtual environment](https://docs.python.org/3/library/venv.html).
 If you're unfamiliar with Python virtual environments, take a look at this [guide](https://packaging.python.org/guides/installing-using-pip-and-virtual-environments/).
-A virtual environment makes it easier to manage different projects and avoid compatibility issues between dependencies.
+A virtual environment makes it easier to manage different projects, and avoid compatibility issues between dependencies.
 
 Start by creating a virtual environment in your project directory:
 
@@ -127,7 +127,7 @@ Your Python environment will find the `main` version of 🤗 Diffusers on the ne
 
 Our library gathers telemetry information during `from_pretrained()` requests.
 This data includes the version of Diffusers and PyTorch/Flax, the requested model or pipeline class,
-and the path to a pre-trained checkpoint if it is hosted on the Hub.
+and the path to a pretrained checkpoint if it is hosted on the Hub.
 This usage data helps us debug issues and prioritize new features.
 Telemetry is only sent when loading models and pipelines from the HuggingFace Hub,
 and is not collected during local usage.
@@ -143,4 +143,4 @@ export DISABLE_TELEMETRY=YES
 On Windows:
 ```bash
 set DISABLE_TELEMETRY=YES
-```
+```

docs/source/en/optimization/habana.mdx

@@ -16,8 +16,8 @@ specific language governing permissions and limitations under the License.
 
 ## Requirements
 
-- Optimum Habana 1.6 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
-- SynapseAI 1.10.
+- Optimum Habana 1.5 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
+- SynapseAI 1.9.
 
 
 ## Inference Pipeline
@@ -41,7 +41,7 @@ pipeline = GaudiStableDiffusionPipeline.from_pretrained(
     scheduler=scheduler,
     use_habana=True,
     use_hpu_graphs=True,
-    gaudi_config="Habana/stable-diffusion-2",
+    gaudi_config="Habana/stable-diffusion",
 )
 ```
 
@@ -62,18 +62,18 @@ For more information, check out Optimum Habana's [documentation](https://hugging
 
 ## Benchmark
 
-Here are the latencies for Habana first-generation Gaudi and Gaudi2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) and [Habana/stable-diffusion-2](https://huggingface.co/Habana/stable-diffusion-2) Gaudi configurations (mixed precision bf16/fp32):
+Here are the latencies for Habana first-generation Gaudi and Gaudi2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) Gaudi configuration (mixed precision bf16/fp32):
 
 - [Stable Diffusion v1.5](https://huggingface.co/runwayml/stable-diffusion-v1-5) (512x512 resolution):
 
 |                        | Latency (batch size = 1) | Throughput (batch size = 8) |
 | ---------------------- |:------------------------:|:---------------------------:|
-| first-generation Gaudi | 3.80s                    | 0.308 images/s              |
-| Gaudi2                 | 1.33s                    | 1.081 images/s              |
+| first-generation Gaudi | 4.22s                    | 0.29 images/s               |
+| Gaudi2                 | 1.70s                    | 0.925 images/s              |
 
 - [Stable Diffusion v2.1](https://huggingface.co/stabilityai/stable-diffusion-2-1) (768x768 resolution):
 
 |                        | Latency (batch size = 1) | Throughput                      |
 | ---------------------- |:------------------------:|:-------------------------------:|
-| first-generation Gaudi | 10.2s                    | 0.108 images/s (batch size = 4) |
-| Gaudi2                 | 3.17s                    | 0.379 images/s (batch size = 8) |
+| first-generation Gaudi | 23.3s                    | 0.045 images/s (batch size = 2) |
+| Gaudi2                 | 7.75s                    | 0.14 images/s (batch size = 5)  |

docs/source/en/quicktour.mdx

@@ -32,9 +32,8 @@ The quicktour is a simplified version of the introductory 🧨 Diffusers [notebo
 
 Before you begin, make sure you have all the necessary libraries installed:
 
-```py
-# uncomment to install the necessary libraries in Colab
-#!pip install --upgrade diffusers accelerate transformers
+```bash
+!pip install --upgrade diffusers accelerate transformers
 ```
 
 - [🤗 Accelerate](https://huggingface.co/docs/accelerate/index) speeds up model loading for inference and training.

docs/source/en/stable_diffusion.mdx

@@ -52,8 +52,6 @@ pipeline = pipeline.to("cuda")
 To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reproducibility):
 
 ```python
-import torch
-
 generator = torch.Generator("cuda").manual_seed(0)
 ```
 

docs/source/en/training/dreambooth.mdx

@@ -12,6 +12,8 @@ specific language governing permissions and limitations under the License.
 
 # DreamBooth
 
+[[open-in-colab]]
+
 [DreamBooth](https://arxiv.org/abs/2208.12242) is a method to personalize text-to-image models like Stable Diffusion given just a few (3-5) images of a subject. It allows the model to generate contextualized images of the subject in different scenes, poses, and views.
 
 ![Dreambooth examples from the project's blog](https://dreambooth.github.io/DreamBooth_files/teaser_static.jpg)
@@ -500,68 +502,9 @@ You may also run inference from any of the [saved training checkpoints](#inferen
 
 ## IF
 
-You can use the lora and full dreambooth scripts to train the text to image [IF model](https://huggingface.co/DeepFloyd/IF-I-XL-v1.0) and the stage II upscaler 
-[IF model](https://huggingface.co/DeepFloyd/IF-II-L-v1.0).
+You can use the lora and full dreambooth scripts to also train the text to image [IF model](https://huggingface.co/DeepFloyd/IF-I-XL-v1.0). A few alternative cli flags are needed due to the model size, the expected input resolution, and the text encoder conventions.
 
-Note that IF has a predicted variance, and our finetuning scripts only train the models predicted error, so for finetuned IF models we switch to a fixed
-variance schedule. The full finetuning scripts will update the scheduler config for the full saved model. However, when loading saved LoRA weights, you
-must also update the pipeline's scheduler config.
-
-```py
-from diffusers import DiffusionPipeline
-
-pipe = DiffusionPipeline.from_pretrained("DeepFloyd/IF-I-XL-v1.0")
-
-pipe.load_lora_weights("<lora weights path>")
-
-# Update scheduler config to fixed variance schedule
-pipe.scheduler = pipe.scheduler.__class__.from_config(pipe.scheduler.config, variance_type="fixed_small")
-```
-
-Additionally, a few alternative cli flags are needed for IF.
-
-`--resolution=64`: IF is a pixel space diffusion model. In order to operate on un-compressed pixels, the input images are of a much smaller resolution. 
-
-`--pre_compute_text_embeddings`: IF uses [T5](https://huggingface.co/docs/transformers/model_doc/t5) for its text encoder. In order to save GPU memory, we pre compute all text embeddings and then de-allocate
-T5.
-
-`--tokenizer_max_length=77`: T5 has a longer default text length, but the default IF encoding procedure uses a smaller number.
-
-`--text_encoder_use_attention_mask`: T5 passes the attention mask to the text encoder.
-
-### Tips and Tricks
-We find LoRA to be sufficient for finetuning the stage I model as the low resolution of the model makes representing finegrained detail hard regardless.
-
-For common and/or not-visually complex object concepts, you can get away with not-finetuning the upscaler. Just be sure to adjust the prompt passed to the
-upscaler to remove the new token from the instance prompt. I.e. if your stage I prompt is "a sks dog", use "a dog" for your stage II prompt.
-
-For finegrained detail like faces that aren't present in the original training set, we find that full finetuning of the stage II upscaler is better than 
-LoRA finetuning stage II.
-
-For finegrained detail like faces, we find that lower learning rates along with larger batch sizes work best.
-
-For stage II, we find that lower learning rates are also needed.
-
-We found experimentally that the DDPM scheduler with the default larger number of denoising steps to sometimes work better than the DPM Solver scheduler
-used in the training scripts.
-
-### Stage II additional validation images
-
-The stage II validation requires images to upscale, we can download a downsized version of the training set:
-
-```py
-from huggingface_hub import snapshot_download
-
-local_dir = "./dog_downsized"
-snapshot_download(
-    "diffusers/dog-example-downsized",
-    local_dir=local_dir,
-    repo_type="dataset",
-    ignore_patterns=".gitattributes",
-)
-```
-
-### IF stage I LoRA Dreambooth
+### LoRA Dreambooth
 This training configuration requires ~28 GB VRAM.
 
 ```sh
@@ -575,7 +518,7 @@ accelerate launch train_dreambooth_lora.py \
   --instance_data_dir=$INSTANCE_DIR \
   --output_dir=$OUTPUT_DIR \
   --instance_prompt="a sks dog" \
-  --resolution=64 \
+  --resolution=64 \ # The input resolution of the IF unet is 64x64
   --train_batch_size=4 \
   --gradient_accumulation_steps=1 \
   --learning_rate=5e-6 \
@@ -584,58 +527,16 @@ accelerate launch train_dreambooth_lora.py \
   --validation_prompt="a sks dog" \
   --validation_epochs=25 \
   --checkpointing_steps=100 \
-  --pre_compute_text_embeddings \
-  --tokenizer_max_length=77 \
-  --text_encoder_use_attention_mask
+  --pre_compute_text_embeddings \ # Pre compute text embeddings to that T5 doesn't have to be kept in memory
+  --tokenizer_max_length=77 \ # IF expects an override of the max token length
+  --text_encoder_use_attention_mask # IF expects attention mask for text embeddings
 ```
 
-### IF stage II LoRA Dreambooth
+### Full Dreambooth
+Due to the size of the optimizer states, we recommend training the full XL IF model with 8bit adam. 
+Using 8bit adam and the rest of the following config, the model can be trained in ~48 GB VRAM.
 
-`--validation_images`: These images are upscaled during validation steps.
-
-`--class_labels_conditioning=timesteps`: Pass additional conditioning to the UNet needed for stage II.
-
-`--learning_rate=1e-6`: Lower learning rate than stage I.
-
-`--resolution=256`: The upscaler expects higher resolution inputs
-
-```sh
-export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
-export INSTANCE_DIR="dog"
-export OUTPUT_DIR="dreambooth_dog_upscale"
-export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
-
-python train_dreambooth_lora.py \
-    --report_to wandb \
-    --pretrained_model_name_or_path=$MODEL_NAME \
-    --instance_data_dir=$INSTANCE_DIR \
-    --output_dir=$OUTPUT_DIR \
-    --instance_prompt="a sks dog" \
-    --resolution=256 \
-    --train_batch_size=4 \
-    --gradient_accumulation_steps=1 \
-    --learning_rate=1e-6 \ 
-    --max_train_steps=2000 \
-    --validation_prompt="a sks dog" \
-    --validation_epochs=100 \
-    --checkpointing_steps=500 \
-    --pre_compute_text_embeddings \
-    --tokenizer_max_length=77 \
-    --text_encoder_use_attention_mask \
-    --validation_images $VALIDATION_IMAGES \
-    --class_labels_conditioning=timesteps
-```
-
-### IF Stage I Full Dreambooth
-`--skip_save_text_encoder`: When training the full model, this will skip saving the entire T5 with the finetuned model. You can still load the pipeline
-with a T5 loaded from the original model.
-
-`use_8bit_adam`: Due to the size of the optimizer states, we recommend training the full XL IF model with 8bit adam. 
-
-`--learning_rate=1e-7`: For full dreambooth, IF requires very low learning rates. With higher learning rates model quality will degrade. Note that it is 
-likely the learning rate can be increased with larger batch sizes.
-
-Using 8bit adam and a batch size of 4, the model can be trained in ~48 GB VRAM.
+For full dreambooth, IF requires very low learning rates. With higher learning rates model quality will degrade.
 
 ```sh
 export MODEL_NAME="DeepFloyd/IF-I-XL-v1.0"
@@ -648,56 +549,17 @@ accelerate launch train_dreambooth.py \
   --instance_data_dir=$INSTANCE_DIR \
   --output_dir=$OUTPUT_DIR \
   --instance_prompt="a photo of sks dog" \
-  --resolution=64 \
+  --resolution=64 \ # The input resolution of the IF unet is 64x64
   --train_batch_size=4 \
   --gradient_accumulation_steps=1 \
   --learning_rate=1e-7 \
   --max_train_steps=150 \
   --validation_prompt "a photo of sks dog" \
   --validation_steps 25 \
-  --text_encoder_use_attention_mask \
-  --tokenizer_max_length 77 \
-  --pre_compute_text_embeddings \
-  --use_8bit_adam \
+  --text_encoder_use_attention_mask \ # IF expects attention mask for text embeddings
+  --tokenizer_max_length 77 \ # IF expects an override of the max token length
+  --pre_compute_text_embeddings \ # Pre compute text embeddings to that T5 doesn't have to be kept in memory
+  --use_8bit_adam \ # 
   --set_grads_to_none \
-  --skip_save_text_encoder \
-  --push_to_hub
-```
-
-### IF Stage II Full Dreambooth
-
-`--learning_rate=5e-6`: With a smaller effective batch size of 4, we found that we required learning rates as low as
-1e-8.
-
-`--resolution=256`: The upscaler expects higher resolution inputs
-
-`--train_batch_size=2` and `--gradient_accumulation_steps=6`: We found that full training of stage II particularly with
-faces required large effective batch sizes.
-
-```sh
-export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
-export INSTANCE_DIR="dog"
-export OUTPUT_DIR="dreambooth_dog_upscale"
-export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
-
-accelerate launch train_dreambooth.py \
-  --report_to wandb \
-  --pretrained_model_name_or_path=$MODEL_NAME \
-  --instance_data_dir=$INSTANCE_DIR \
-  --output_dir=$OUTPUT_DIR \
-  --instance_prompt="a sks dog" \
-  --resolution=256 \
-  --train_batch_size=2 \
-  --gradient_accumulation_steps=6 \
-  --learning_rate=5e-6 \
-  --max_train_steps=2000 \
-  --validation_prompt="a sks dog" \
-  --validation_steps=150 \
-  --checkpointing_steps=500 \
-  --pre_compute_text_embeddings \
-  --tokenizer_max_length=77 \
-  --text_encoder_use_attention_mask \
-  --validation_images $VALIDATION_IMAGES \
-  --class_labels_conditioning timesteps \
-  --push_to_hub
-```
+  --skip_save_text_encoder # do not save the full T5 text encoder with the model
+```

docs/source/en/training/instructpix2pix.mdx

@@ -207,5 +207,3 @@ speed and quality during performance:
 
 Particularly, `image_guidance_scale` and `guidance_scale` can have a profound impact
 on the generated ("edited") image (see [here](https://twitter.com/RisingSayak/status/1628392199196151808?s=20) for an example).
-
-If you're looking for some interesting ways to use the InstructPix2Pix training methodology, we welcome you to check out this blog post: [Instruction-tuning Stable Diffusion with InstructPix2Pix](https://huggingface.co/blog/instruction-tuning-sd). 

docs/source/en/training/lora.mdx

@@ -12,6 +12,8 @@ specific language governing permissions and limitations under the License.
 
 # Low-Rank Adaptation of Large Language Models (LoRA)
 
+[[open-in-colab]]
+
 <Tip warning={true}>
 
 Currently, LoRA is only supported for the attention layers of the [`UNet2DConditionalModel`]. We also 
@@ -258,14 +260,6 @@ pipe.load_lora_weights(lora_model_id)
 image = pipe("A picture of a sks dog in a bucket", num_inference_steps=25).images[0]
 ```
 
-<Tip>
-
-If your LoRA parameters involve the UNet as well as the Text Encoder, then passing
-`cross_attention_kwargs={"scale": 0.5}` will apply the `scale` value to both the UNet 
-and the Text Encoder. 
-
-</Tip>
-
 Note that the use of [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] is preferred to [`~diffusers.loaders.UNet2DConditionLoadersMixin.load_attn_procs`] for loading LoRA parameters. This is because
 [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] can handle the following situations:
 
@@ -278,75 +272,4 @@ Note that the use of [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] is
 * LoRA parameters that have separate identifiers for the UNet and the text encoder such as: [`"sayakpaul/dreambooth"`](https://huggingface.co/sayakpaul/dreambooth).
 
 **Note** that it is possible to provide a local directory path to [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] as well as [`~diffusers.loaders.UNet2DConditionLoadersMixin.load_attn_procs`]. To know about the supported inputs,
-refer to the respective docstrings.
-
-## Supporting A1111 themed LoRA checkpoints from Diffusers
-
-To provide seamless interoperability with A1111 to our users, we support loading A1111 formatted
-LoRA checkpoints using [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] in a limited capacity.
-In this section, we explain how to load an A1111 formatted LoRA checkpoint from [CivitAI](https://civitai.com/)
-in Diffusers and perform inference with it. 
-
-First, download a checkpoint. We'll use
-[this one](https://civitai.com/models/13239/light-and-shadow) for demonstration purposes. 
-
-```bash
-wget https://civitai.com/api/download/models/15603 -O light_and_shadow.safetensors
-```
-
-Next, we initialize a [`~DiffusionPipeline`]:
-
-```python 
-import torch
-
-from diffusers import StableDiffusionPipeline, DPMSolverMultistepScheduler
-
-pipeline = StableDiffusionPipeline.from_pretrained(
-    "gsdf/Counterfeit-V2.5", torch_dtype=torch.float16, safety_checker=None
-).to("cuda")
-pipeline.scheduler = DPMSolverMultistepScheduler.from_config(
-    pipeline.scheduler.config, use_karras_sigmas=True
-)
-```
-
-We then load the checkpoint downloaded from CivitAI: 
-
-```python 
-pipeline.load_lora_weights(".", weight_name="light_and_shadow.safetensors")
-```
-
-<Tip warning={true}>
-
-If you're loading a checkpoint in the `safetensors` format, please ensure you have `safetensors` installed.
-
-</Tip>
-
-And then it's time for running inference: 
-
-```python 
-prompt = "masterpiece, best quality, 1girl, at dusk"
-negative_prompt = ("(low quality, worst quality:1.4), (bad anatomy), (inaccurate limb:1.2), "
-                   "bad composition, inaccurate eyes, extra digit, fewer digits, (extra arms:1.2), large breasts")
-
-images = pipeline(prompt=prompt, 
-    negative_prompt=negative_prompt, 
-    width=512, 
-    height=768, 
-    num_inference_steps=15, 
-    num_images_per_prompt=4,
-    generator=torch.manual_seed(0)
-).images
-```
-
-Below is a comparison between the LoRA and the non-LoRA results:
-
-![lora_non_lora](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lora_non_lora_comparison.png)
-
-You have a similar checkpoint stored on the Hugging Face Hub, you can load it
-directly with [`~diffusers.loaders.LoraLoaderMixin.load_lora_weights`] like so: 
-
-```python 
-lora_model_id = "sayakpaul/civitai-light-shadow-lora"
-lora_filename = "light_and_shadow.safetensors"
-pipeline.load_lora_weights(lora_model_id, weight_name=lora_filename)
-```
+refer to the respective docstrings.

docs/source/en/training/text_inversion.mdx

@@ -14,6 +14,8 @@ specific language governing permissions and limitations under the License.
 
 # Textual Inversion
 
+[[open-in-colab]]
+
 [Textual Inversion](https://arxiv.org/abs/2208.01618) is a technique for capturing novel concepts from a small number of example images. While the technique was originally demonstrated with a [latent diffusion model](https://github.com/CompVis/latent-diffusion), it has since been applied to other model variants like [Stable Diffusion](https://huggingface.co/docs/diffusers/main/en/conceptual/stable_diffusion). The learned concepts can be used to better control the images generated from text-to-image pipelines. It learns new "words" in the text encoder's embedding space, which are used within text prompts for personalized image generation.
 
 ![Textual Inversion example](https://textual-inversion.github.io/static/images/editing/colorful_teapot.JPG)

docs/source/en/tutorials/basic_training.mdx

@@ -26,9 +26,8 @@ This tutorial will teach you how to train a [`UNet2DModel`] from scratch on a su
 
 Before you begin, make sure you have 🤗 Datasets installed to load and preprocess image datasets, and 🤗 Accelerate, to simplify training on any number of GPUs. The following command will also install [TensorBoard](https://www.tensorflow.org/tensorboard) to visualize training metrics (you can also use [Weights & Biases](https://docs.wandb.ai/) to track your training).
 
-```py
-# uncomment to install the necessary libraries in Colab
-#!pip install diffusers[training]
+```bash
+!pip install diffusers[training]
 ```
 
 We encourage you to share your model with the community, and in order to do that, you'll need to login to your Hugging Face account (create one [here](https://hf.co/join) if you don't already have one!). You can login from a notebook and enter your token when prompted:
@@ -313,7 +312,7 @@ Now you can wrap all these components together in a training loop with 🤗 Acce
 ...         mixed_precision=config.mixed_precision,
 ...         gradient_accumulation_steps=config.gradient_accumulation_steps,
 ...         log_with="tensorboard",
-...         project_dir=os.path.join(config.output_dir, "logs"),
+...         logging_dir=os.path.join(config.output_dir, "logs"),
 ...     )
 ...     if accelerator.is_main_process:
 ...         if config.push_to_hub:

docs/source/en/using-diffusers/control_brightness.mdx

@@ -1,45 +0,0 @@
-# Control image brightness
-
-The Stable Diffusion pipeline is mediocre at generating images that are either very bright or dark as explained in the [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) paper. The solutions proposed in the paper are currently implemented in the [`DDIMScheduler`] which you can use to improve the lighting in your images.
-
-<Tip>
-
-💡 Take a look at the paper linked above for more details about the proposed solutions!
-
-</Tip>
-
-One of the solutions is to train a model with *v prediction* and *v loss*. Add the following flag to the [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py) scripts to enable `v_prediction`:
-
-```bash
---prediction_type="v_prediction"
-```
-
-For example, let's use the [`ptx0/pseudo-journey-v2`](https://huggingface.co/ptx0/pseudo-journey-v2) checkpoint which has been finetuned with `v_prediction`.
-
-Next, configure the following parameters in the [`DDIMScheduler`]:
-
-1. `rescale_betas_zero_snr=True`, rescales the noise schedule to zero terminal signal-to-noise ratio (SNR)
-2. `timestep_spacing="trailing"`, starts sampling from the last timestep
-
-```py
->>> from diffusers import DiffusionPipeline, DDIMScheduler
-
->>> pipeline = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2")
-# switch the scheduler in the pipeline to use the DDIMScheduler
-
->>> pipeline.scheduler = DDIMScheduler.from_config(
-...     pipeline.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
-... )
->>> pipeline.to("cuda")
-```
-
-Finally, in your call to the pipeline, set `guidance_rescale` to prevent overexposure:
-
-```py
-prompt = "A lion in galaxies, spirals, nebulae, stars, smoke, iridescent, intricate detail, octane render, 8k"
-image = pipeline(prompt, guidance_rescale=0.7).images[0]
-```
-
-<div class="flex justify-center">
-    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/zero_snr.png"/>
-</div>

docs/source/en/using-diffusers/custom_pipeline_examples.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Community pipelines
 
-[[open-in-colab]]
-
 > **For more information about community pipelines, please have a look at [this issue](https://github.com/huggingface/diffusers/issues/841).**
 
 **Community** examples consist of both inference and training examples that have been added by the community.

docs/source/en/using-diffusers/custom_pipeline_overview.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Load community pipelines
 
-[[open-in-colab]]
-
 Community pipelines are any [`DiffusionPipeline`] class that are different from the original implementation as specified in their paper (for example, the [`StableDiffusionControlNetPipeline`] corresponds to the [Text-to-Image Generation with ControlNet Conditioning](https://arxiv.org/abs/2302.05543) paper). They provide additional functionality or extend the original implementation of a pipeline.
 
 There are many cool community pipelines like [Speech to Image](https://github.com/huggingface/diffusers/tree/main/examples/community#speech-to-image) or [Composable Stable Diffusion](https://github.com/huggingface/diffusers/tree/main/examples/community#composable-stable-diffusion), and you can find all the official community pipelines [here](https://github.com/huggingface/diffusers/tree/main/examples/community).

docs/source/en/using-diffusers/img2img.mdx

@@ -18,9 +18,8 @@ The [`StableDiffusionImg2ImgPipeline`] lets you pass a text prompt and an initia
 
 Before you begin, make sure you have all the necessary libraries installed:
 
-```py
-# uncomment to install the necessary libraries in Colab
-#!pip install diffusers transformers ftfy accelerate
+```bash
+!pip install diffusers transformers ftfy accelerate
 ```
 
 Get started by creating a [`StableDiffusionImg2ImgPipeline`] with a pretrained Stable Diffusion model like [`nitrosocke/Ghibli-Diffusion`](https://huggingface.co/nitrosocke/Ghibli-Diffusion).

docs/source/en/using-diffusers/loading.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Load pipelines, models, and schedulers
 
-[[open-in-colab]]
-
 Having an easy way to use a diffusion system for inference is essential to 🧨 Diffusers. Diffusion systems often consist of multiple components like parameterized models, tokenizers, and schedulers that interact in complex ways. That is why we designed the [`DiffusionPipeline`] to wrap the complexity of the entire diffusion system into an easy-to-use API, while remaining flexible enough to be adapted for other use cases, such as loading each component individually as building blocks to assemble your own diffusion system.
 
 Everything you need for inference or training is accessible with the `from_pretrained()` method.
@@ -174,7 +172,7 @@ A checkpoint variant is usually a checkpoint where it's weights are:
 
 </Tip>
 
-Otherwise, a variant is **identical** to the original checkpoint. They have exactly the same serialization format (like [Safetensors](./using_safetensors)), model structure, and weights have identical tensor shapes.
+Otherwise, a variant is **identical** to the original checkpoint. They have exactly the same serialization format (like [Safetensors](./using-diffusers/using_safetensors)), model structure, and weights have identical tensor shapes.
 
 | **checkpoint type** | **weight name**                     | **argument for loading weights** |
 |---------------------|-------------------------------------|----------------------------------|
@@ -190,7 +188,6 @@ There are two important arguments to know for loading variants:
 
 ```python
 from diffusers import DiffusionPipeline
-import torch
 
 # load fp16 variant
 stable_diffusion = DiffusionPipeline.from_pretrained(

docs/source/en/using-diffusers/other-formats.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Load different Stable Diffusion formats
 
-[[open-in-colab]]
-
 Stable Diffusion models are available in different formats depending on the framework they're trained and saved with, and where you download them from. Converting these formats for use in 🤗 Diffusers allows you to use all the features supported by the library, such as [using different schedulers](schedulers) for inference, [building your custom pipeline](write_own_pipeline), and a variety of techniques and methods for [optimizing inference speed](./optimization/opt_overview).
 
 <Tip>
@@ -26,7 +24,7 @@ This guide will show you how to convert other Stable Diffusion formats to be com
 
 ## PyTorch .ckpt
 
-The checkpoint - or `.ckpt` - format is commonly used to store and save models. The `.ckpt` file contains the entire model and is typically several GBs in size. While you can load and use a `.ckpt` file directly with the [`~StableDiffusionPipeline.from_single_file`] method, it is generally better to convert the `.ckpt` file to 🤗 Diffusers so both formats are available.
+The checkpoint - or `.ckpt` - format is commonly used to store and save models. The `.ckpt` file contains the entire model and is typically several GBs in size. While you can load and use a `.ckpt` file directly with the [`~StableDiffusionPipeline.from_ckpt`] method, it is generally better to convert the `.ckpt` file to 🤗 Diffusers so both formats are available.
 
 There are two options for converting a `.ckpt` file; use a Space to convert the checkpoint or convert the `.ckpt` file with a script.
 
@@ -125,70 +123,4 @@ pipeline.to("cuda")
 placeholder_token = "<my-funny-cat-token>"
 prompt = f"two {placeholder_token} getting married, photorealistic, high quality"
 image = pipeline(prompt, num_inference_steps=50).images[0]
-```
-
-## A1111 LoRA files
-
-[Automatic1111](https://github.com/AUTOMATIC1111/stable-diffusion-webui) (A1111) is a popular web UI for Stable Diffusion that supports model sharing platforms like [Civitai](https://civitai.com/). Models trained with the Low-Rank Adaptation (LoRA) technique are especially popular because they're fast to train and have a much smaller file size than a fully finetuned model. 🤗 Diffusers supports loading A1111 LoRA checkpoints with [`~loaders.LoraLoaderMixin.load_lora_weights`]:
-
-```py
-from diffusers import DiffusionPipeline, UniPCMultistepScheduler
-import torch
-
-pipeline = DiffusionPipeline.from_pretrained(
-    "andite/anything-v4.0", torch_dtype=torch.float16, safety_checker=None
-).to("cuda")
-pipeline.scheduler = UniPCMultistepScheduler.from_config(pipeline.scheduler.config)
-```
-
-Download a LoRA checkpoint from Civitai; this example uses the [Howls Moving Castle,Interior/Scenery LoRA (Ghibli Stlye)](https://civitai.com/models/14605?modelVersionId=19998) checkpoint, but feel free to try out any LoRA checkpoint!
-
-```py
-# uncomment to download the safetensor weights
-#!wget https://civitai.com/api/download/models/19998 -O howls_moving_castle.safetensors
-```
-
-Load the LoRA checkpoint into the pipeline with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method:
-
-```py
-pipeline.load_lora_weights(".", weight_name="howls_moving_castle.safetensors")
-```
-
-Now you can use the pipeline to generate images:
-
-```py
-prompt = "masterpiece, illustration, ultra-detailed, cityscape, san francisco, golden gate bridge, california, bay area, in the snow, beautiful detailed starry sky"
-negative_prompt = "lowres, cropped, worst quality, low quality, normal quality, artifacts, signature, watermark, username, blurry, more than one bridge, bad architecture"
-
-images = pipeline(
-    prompt=prompt,
-    negative_prompt=negative_prompt,
-    width=512,
-    height=512,
-    num_inference_steps=25,
-    num_images_per_prompt=4,
-    generator=torch.manual_seed(0),
-).images
-```
-
-Finally, create a helper function to display the images:
-
-```py
-from PIL import Image
-
-
-def image_grid(imgs, rows=2, cols=2):
-    w, h = imgs[0].size
-    grid = Image.new("RGB", size=(cols * w, rows * h))
-
-    for i, img in enumerate(imgs):
-        grid.paste(img, box=(i % cols * w, i // cols * h))
-    return grid
-
-
-image_grid(images)
-```
-
-<div class="flex justify-center">
-    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/a1111-lora-sf.png"/>
-</div>
+```

docs/source/en/using-diffusers/reproducibility.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Create reproducible pipelines
 
-[[open-in-colab]]
-
 Reproducibility is important for testing, replicating results, and can even be used to [improve image quality](reusing_seeds). However, the randomness in diffusion models is a desired property because it allows the pipeline to generate different images every time it is run. While you can't expect to get the exact same results across platforms, you can expect results to be reproducible across releases and platforms within a certain tolerance range. Even then, tolerance varies depending on the diffusion pipeline and checkpoint.
 
 This is why it's important to understand how to control sources of randomness in diffusion models or use deterministic algorithms.
@@ -113,7 +111,7 @@ print(np.abs(image).sum())
 
 The result is not the same even though you're using an identical seed because the GPU uses a different random number generator than the CPU.
 
-To circumvent this problem, 🧨 Diffusers has a [`~diffusers.utils.randn_tensor`] function for creating random noise on the CPU, and then moving the tensor to a GPU if necessary. The `randn_tensor` function is used everywhere inside the pipeline, allowing the user to **always** pass a CPU `Generator` even if the pipeline is run on a GPU. 
+To circumvent this problem, 🧨 Diffusers has a [`randn_tensor`](#diffusers.utils.randn_tensor) function for creating random noise on the CPU, and then moving the tensor to a GPU if necessary. The `randn_tensor` function is used everywhere inside the pipeline, allowing the user to **always** pass a CPU `Generator` even if the pipeline is run on a GPU. 
 
 You'll see the results are much closer now!
 
@@ -149,6 +147,9 @@ susceptible to precision error propagation. Don't expect similar results across
 different GPU hardware or PyTorch versions. In this case, you'll need to run 
 exactly the same hardware and PyTorch version for full reproducibility.
 
+### randn_tensor
+[[autodoc]] diffusers.utils.randn_tensor
+
 ## Deterministic algorithms
 
 You can also configure PyTorch to use deterministic algorithms to create a reproducible pipeline. However, you should be aware that deterministic algorithms may be slower than nondeterministic ones and you may observe a decrease in performance. But if reproducibility is important to you, then this is the way to go!

docs/source/en/using-diffusers/reusing_seeds.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Improve image quality with deterministic generation
 
-[[open-in-colab]]
-
 A common way to improve the quality of generated images is with *deterministic batch generation*, generate a batch of images and select one image to improve with a more detailed prompt in a second round of inference. The key is to pass a list of [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html#generator)'s to the pipeline for batched image generation, and tie each `Generator` to a seed so you can reuse it for an image.
 
 Let's use [`runwayml/stable-diffusion-v1-5`](runwayml/stable-diffusion-v1-5) for example, and generate several versions of the following prompt:

docs/source/en/using-diffusers/schedulers.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Schedulers
 
-[[open-in-colab]]
-
 Diffusion pipelines are inherently a collection of diffusion models and schedulers that are partly independent from each other. This means that one is able to switch out parts of the pipeline to better customize 
 a pipeline to one's use case. The best example of this is the [Schedulers](../api/schedulers/overview.mdx).
 

docs/source/en/using-diffusers/stable_diffusion_jax_how_to.mdx

@@ -14,10 +14,9 @@ Note that JAX is not exclusive to TPUs, but it shines on that hardware because e
 
 First make sure diffusers is installed.
 
-```py
-# uncomment to install the necessary libraries in Colab
-#!pip install jax==0.3.25 jaxlib==0.3.25 flax transformers ftfy
-#!pip install diffusers
+```bash
+!pip install jax==0.3.25 jaxlib==0.3.25 flax transformers ftfy
+!pip install diffusers
 ```
 
 ```python

docs/source/en/using-diffusers/using_safetensors.mdx

@@ -1,14 +1,11 @@
 # Load safetensors
 
-[[open-in-colab]]
-
 [safetensors](https://github.com/huggingface/safetensors) is a safe and fast file format for storing and loading tensors. Typically, PyTorch model weights are saved or *pickled* into a `.bin` file with Python's [`pickle`](https://docs.python.org/3/library/pickle.html) utility. However, `pickle` is not secure and pickled files may contain malicious code that can be executed. safetensors is a secure alternative to `pickle`, making it ideal for sharing model weights.
 
 This guide will show you how you load `.safetensor` files, and how to convert Stable Diffusion model weights stored in other formats to `.safetensor`. Before you start, make sure you have safetensors installed:
 
-```py
-# uncomment to install the necessary libraries in Colab
-#!pip install safetensors
+```bash
+!pip install safetensors
 ```
 
 If you look at the [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5/tree/main) repository, you'll see weights inside the `text_encoder`, `unet` and `vae` subfolders are stored in the `.safetensors` format. By default, 🤗 Diffusers automatically loads these `.safetensors` files from their subfolders if they're available in the model repository.
@@ -21,12 +18,12 @@ from diffusers import DiffusionPipeline
 pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True)
 ```
 
-However, model weights are not necessarily stored in separate subfolders like in the example above. Sometimes, all the weights are stored in a single `.safetensors` file. In this case, if the weights are Stable Diffusion weights, you can load the file directly with the [`~diffusers.loaders.FromSingleFileMixin.from_single_file`] method:
+However, model weights are not necessarily stored in separate subfolders like in the example above. Sometimes, all the weights are stored in a single `.safetensors` file. In this case, if the weights are Stable Diffusion weights, you can load the file directly with the [`~diffusers.loaders.FromCkptMixin.from_ckpt`] method:
 
 ```py
 from diffusers import StableDiffusionPipeline
 
-pipeline = StableDiffusionPipeline.from_single_file(
+pipeline = StableDiffusionPipeline.from_ckpt(
     "https://huggingface.co/WarriorMama777/OrangeMixs/blob/main/Models/AbyssOrangeMix/AbyssOrangeMix.safetensors"
 )
 ```

docs/source/en/using-diffusers/weighted_prompts.mdx

@@ -12,8 +12,6 @@ specific language governing permissions and limitations under the License.
 
 # Weighting prompts
 
-[[open-in-colab]]
-
 Text-guided diffusion models generate images based on a given text prompt. The text prompt
 can include multiple concepts that the model should generate and it's often desirable to weight
 certain parts of the prompt more or less. 
@@ -96,15 +94,5 @@ a try!
 If your favorite pipeline does not have a `prompt_embeds` input, please make sure to open an issue, the 
 diffusers team tries to be as responsive as possible.
 
-Compel 1.1.6 adds a utility class to simplify using textual inversions.  Instantiate a `DiffusersTextualInversionManager` and pass it to Compel init:
-
-```
-textual_inversion_manager = DiffusersTextualInversionManager(pipe)
-compel = Compel(
-    tokenizer=pipe.tokenizer,
-    text_encoder=pipe.text_encoder,
-    textual_inversion_manager=textual_inversion_manager)
-```
-
 Also, please check out the documentation of the [compel](https://github.com/damian0815/compel) library for 
 more information.

docs/source/en/using-diffusers/write_own_pipeline.mdx

@@ -42,63 +42,63 @@ To recreate the pipeline with the model and scheduler separately, let's write ou
 
 1. Load the model and scheduler:
 
-```py
->>> from diffusers import DDPMScheduler, UNet2DModel
+    ```py
+    >>> from diffusers import DDPMScheduler, UNet2DModel
 
->>> scheduler = DDPMScheduler.from_pretrained("google/ddpm-cat-256")
->>> model = UNet2DModel.from_pretrained("google/ddpm-cat-256").to("cuda")
-```
+    >>> scheduler = DDPMScheduler.from_pretrained("google/ddpm-cat-256")
+    >>> model = UNet2DModel.from_pretrained("google/ddpm-cat-256").to("cuda")
+    ```
 
 2. Set the number of timesteps to run the denoising process for:
 
-```py
->>> scheduler.set_timesteps(50)
-```
+    ```py
+    >>> scheduler.set_timesteps(50)
+    ```
 
 3. Setting the scheduler timesteps creates a tensor with evenly spaced elements in it, 50 in this example. Each element corresponds to a timestep at which the model denoises an image. When you create the denoising loop later, you'll iterate over this tensor to denoise an image:
 
-```py
->>> scheduler.timesteps
-tensor([980, 960, 940, 920, 900, 880, 860, 840, 820, 800, 780, 760, 740, 720,
-    700, 680, 660, 640, 620, 600, 580, 560, 540, 520, 500, 480, 460, 440,
-    420, 400, 380, 360, 340, 320, 300, 280, 260, 240, 220, 200, 180, 160,
-    140, 120, 100,  80,  60,  40,  20,   0])
-```
+    ```py
+    >>> scheduler.timesteps
+    tensor([980, 960, 940, 920, 900, 880, 860, 840, 820, 800, 780, 760, 740, 720,
+        700, 680, 660, 640, 620, 600, 580, 560, 540, 520, 500, 480, 460, 440,
+        420, 400, 380, 360, 340, 320, 300, 280, 260, 240, 220, 200, 180, 160,
+        140, 120, 100,  80,  60,  40,  20,   0])
+    ```
 
 4. Create some random noise with the same shape as the desired output:
 
-```py
->>> import torch
+    ```py
+    >>> import torch
 
->>> sample_size = model.config.sample_size
->>> noise = torch.randn((1, 3, sample_size, sample_size)).to("cuda")
-```
+    >>> sample_size = model.config.sample_size
+    >>> noise = torch.randn((1, 3, sample_size, sample_size)).to("cuda")
+    ```
 
-5. Now write a loop to iterate over the timesteps. At each timestep, the model does a [`UNet2DModel.forward`] pass and returns the noisy residual. The scheduler's [`~DDPMScheduler.step`] method takes the noisy residual, timestep, and input and it predicts the image at the previous timestep. This output becomes the next input to the model in the denoising loop, and it'll repeat until it reaches the end of the `timesteps` array.
+4. Now write a loop to iterate over the timesteps. At each timestep, the model does a [`UNet2DModel.forward`] pass and returns the noisy residual. The scheduler's [`~DDPMScheduler.step`] method takes the noisy residual, timestep, and input and it predicts the image at the previous timestep. This output becomes the next input to the model in the denoising loop, and it'll repeat until it reaches the end of the `timesteps` array.
 
-```py
->>> input = noise
+    ```py
+    >>> input = noise
 
->>> for t in scheduler.timesteps:
-...     with torch.no_grad():
-...         noisy_residual = model(input, t).sample
-...     previous_noisy_sample = scheduler.step(noisy_residual, t, input).prev_sample
-...     input = previous_noisy_sample
-```
+    >>> for t in scheduler.timesteps:
+    ...     with torch.no_grad():
+    ...         noisy_residual = model(input, t).sample
+    ...     previous_noisy_sample = scheduler.step(noisy_residual, t, input).prev_sample
+    ...     input = previous_noisy_sample
+    ```
 
-This is the entire denoising process, and you can use this same pattern to write any diffusion system.
+    This is the entire denoising process, and you can use this same pattern to write any diffusion system.
 
-6. The last step is to convert the denoised output into an image:
+5. The last step is to convert the denoised output into an image:
 
-```py
->>> from PIL import Image
->>> import numpy as np
+    ```py
+    >>> from PIL import Image
+    >>> import numpy as np
 
->>> image = (input / 2 + 0.5).clamp(0, 1)
->>> image = image.cpu().permute(0, 2, 3, 1).numpy()[0]
->>> image = Image.fromarray((image * 255).round().astype("uint8"))
->>> image
-```
+    >>> image = (input / 2 + 0.5).clamp(0, 1)
+    >>> image = image.cpu().permute(0, 2, 3, 1).numpy()[0]
+    >>> image = Image.fromarray((image * 255).round().astype("uint8"))
+    >>> image
+    ```
 
 In the next section, you'll put your skills to the test and breakdown the more complex Stable Diffusion pipeline. The steps are more or less the same. You'll initialize the necessary components, and set the number of timesteps to create a `timestep` array. The `timestep` array is used in the denoising loop, and for each element in this array, the model predicts a less noisy image. The denoising loop iterates over the `timestep`'s, and at each timestep, it outputs a noisy residual and the scheduler uses it to predict a less noisy image at the previous timestep. This process is repeated until you reach the end of the `timestep` array.
 
@@ -286,5 +286,5 @@ This is really what 🧨 Diffusers is designed for: to make it intuitive and eas
 
 For your next steps, feel free to:
 
-* Learn how to [build and contribute a pipeline](contribute_pipeline) to 🧨 Diffusers. We can't wait and see what you'll come up with!
-* Explore [existing pipelines](../api/pipelines/overview) in the library, and see if you can deconstruct and build a pipeline from scratch using the models and schedulers separately.
+* Learn how to [build and contribute a pipeline](using-diffusers/#contribute_pipeline) to 🧨 Diffusers. We can't wait and see what you'll come up with!
+* Explore [existing pipelines](./api/pipelines/overview) in the library, and see if you can deconstruct and build a pipeline from scratch using the models and schedulers separately.

docs/source/ko/_toctree.yml

@@ -45,4 +45,4 @@
     title: MPS
   - local: optimization/habana
     title: Habana Gaudi
-  title: 최적화/특수 하드웨어
+  title: 최적화/특수 하드웨어

docs/source/zh/_toctree.yml

@@ -3,6 +3,272 @@
     title: 🧨 Diffusers
   - local: quicktour
     title: 快速入门
+  - local: stable_diffusion
+    title: Effective and efficient diffusion
   - local: installation
     title: 安装
   title: 开始
+- sections:
+  - local: tutorials/tutorial_overview
+    title: Overview
+  - local: using-diffusers/write_own_pipeline
+    title: Understanding models and schedulers
+  - local: tutorials/basic_training
+    title: Train a diffusion model
+  title: Tutorials
+- sections:
+  - sections:
+    - local: using-diffusers/loading_overview
+      title: Overview
+    - local: using-diffusers/loading
+      title: Load pipelines, models, and schedulers
+    - local: using-diffusers/schedulers
+      title: Load and compare different schedulers
+    - local: using-diffusers/custom_pipeline_overview
+      title: Load community pipelines
+    - local: using-diffusers/kerascv
+      title: Load KerasCV Stable Diffusion checkpoints
+    title: Loading & Hub
+  - sections:
+    - local: using-diffusers/pipeline_overview
+      title: Overview
+    - local: using-diffusers/unconditional_image_generation
+      title: Unconditional image generation
+    - local: using-diffusers/conditional_image_generation
+      title: Text-to-image generation
+    - local: using-diffusers/img2img
+      title: Text-guided image-to-image
+    - local: using-diffusers/inpaint
+      title: Text-guided image-inpainting
+    - local: using-diffusers/depth2img
+      title: Text-guided depth-to-image
+    - local: using-diffusers/reusing_seeds
+      title: Improve image quality with deterministic generation
+    - local: using-diffusers/reproducibility
+      title: Create reproducible pipelines
+    - local: using-diffusers/custom_pipeline_examples
+      title: Community pipelines
+    - local: using-diffusers/contribute_pipeline
+      title: How to contribute a community pipeline
+    - local: using-diffusers/using_safetensors
+      title: Using safetensors
+    - local: using-diffusers/stable_diffusion_jax_how_to
+      title: Stable Diffusion in JAX/Flax
+    - local: using-diffusers/weighted_prompts
+      title: Weighting Prompts
+    title: Pipelines for Inference
+  - sections:
+    - local: training/overview
+      title: Overview
+    - local: training/unconditional_training
+      title: Unconditional image generation
+    - local: training/text_inversion
+      title: Textual Inversion
+    - local: training/dreambooth
+      title: DreamBooth
+    - local: training/text2image
+      title: Text-to-image
+    - local: training/lora
+      title: Low-Rank Adaptation of Large Language Models (LoRA)
+    - local: training/controlnet
+      title: ControlNet
+    - local: training/instructpix2pix
+      title: InstructPix2Pix Training
+    - local: training/custom_diffusion
+      title: Custom Diffusion
+    title: Training
+  - sections:
+    - local: using-diffusers/rl
+      title: Reinforcement Learning
+    - local: using-diffusers/audio
+      title: Audio
+    - local: using-diffusers/other-modalities
+      title: Other Modalities
+    title: Taking Diffusers Beyond Images
+  title: Using Diffusers
+- sections:
+  - local: optimization/opt_overview
+    title: Overview
+  - local: optimization/fp16
+    title: Memory and Speed
+  - local: optimization/torch2.0
+    title: Torch2.0 support
+  - local: optimization/xformers
+    title: xFormers
+  - local: optimization/onnx
+    title: ONNX
+  - local: optimization/open_vino
+    title: OpenVINO
+  - local: optimization/coreml
+    title: Core ML
+  - local: optimization/mps
+    title: MPS
+  - local: optimization/habana
+    title: Habana Gaudi
+  - local: optimization/tome
+    title: Token Merging
+  title: Optimization/Special Hardware
+- sections:
+  - local: conceptual/philosophy
+    title: Philosophy
+  - local: using-diffusers/controlling_generation
+    title: Controlled generation
+  - local: conceptual/contribution
+    title: How to contribute?
+  - local: conceptual/ethical_guidelines
+    title: Diffusers' Ethical Guidelines
+  - local: conceptual/evaluation
+    title: Evaluating Diffusion Models
+  title: Conceptual Guides
+- sections:
+  - sections:
+    - local: api/models
+      title: Models
+    - local: api/diffusion_pipeline
+      title: Diffusion Pipeline
+    - local: api/logging
+      title: Logging
+    - local: api/configuration
+      title: Configuration
+    - local: api/outputs
+      title: Outputs
+    - local: api/loaders
+      title: Loaders
+    title: Main Classes
+  - sections:
+    - local: api/pipelines/overview
+      title: Overview
+    - local: api/pipelines/alt_diffusion
+      title: AltDiffusion
+    - local: api/pipelines/audio_diffusion
+      title: Audio Diffusion
+    - local: api/pipelines/audioldm
+      title: AudioLDM
+    - local: api/pipelines/cycle_diffusion
+      title: Cycle Diffusion
+    - local: api/pipelines/dance_diffusion
+      title: Dance Diffusion
+    - local: api/pipelines/ddim
+      title: DDIM
+    - local: api/pipelines/ddpm
+      title: DDPM
+    - local: api/pipelines/dit
+      title: DiT
+    - local: api/pipelines/if
+      title: IF
+    - local: api/pipelines/latent_diffusion
+      title: Latent Diffusion
+    - local: api/pipelines/paint_by_example
+      title: PaintByExample
+    - local: api/pipelines/pndm
+      title: PNDM
+    - local: api/pipelines/repaint
+      title: RePaint
+    - local: api/pipelines/stable_diffusion_safe
+      title: Safe Stable Diffusion
+    - local: api/pipelines/score_sde_ve
+      title: Score SDE VE
+    - local: api/pipelines/semantic_stable_diffusion
+      title: Semantic Guidance
+    - local: api/pipelines/spectrogram_diffusion
+      title: "Spectrogram Diffusion"
+    - sections:
+      - local: api/pipelines/stable_diffusion/overview
+        title: Overview
+      - local: api/pipelines/stable_diffusion/text2img
+        title: Text-to-Image
+      - local: api/pipelines/stable_diffusion/img2img
+        title: Image-to-Image
+      - local: api/pipelines/stable_diffusion/inpaint
+        title: Inpaint
+      - local: api/pipelines/stable_diffusion/depth2img
+        title: Depth-to-Image
+      - local: api/pipelines/stable_diffusion/image_variation
+        title: Image-Variation
+      - local: api/pipelines/stable_diffusion/upscale
+        title: Super-Resolution
+      - local: api/pipelines/stable_diffusion/latent_upscale
+        title: Stable-Diffusion-Latent-Upscaler
+      - local: api/pipelines/stable_diffusion/pix2pix
+        title: InstructPix2Pix
+      - local: api/pipelines/stable_diffusion/attend_and_excite
+        title: Attend and Excite
+      - local: api/pipelines/stable_diffusion/pix2pix_zero
+        title: Pix2Pix Zero
+      - local: api/pipelines/stable_diffusion/self_attention_guidance
+        title: Self-Attention Guidance
+      - local: api/pipelines/stable_diffusion/panorama
+        title: MultiDiffusion Panorama
+      - local: api/pipelines/stable_diffusion/controlnet
+        title: Text-to-Image Generation with ControlNet Conditioning
+      - local: api/pipelines/stable_diffusion/model_editing
+        title: Text-to-Image Model Editing
+      title: Stable Diffusion
+    - local: api/pipelines/stable_diffusion_2
+      title: Stable Diffusion 2
+    - local: api/pipelines/stable_unclip
+      title: Stable unCLIP
+    - local: api/pipelines/stochastic_karras_ve
+      title: Stochastic Karras VE
+    - local: api/pipelines/text_to_video
+      title: Text-to-Video
+    - local: api/pipelines/text_to_video_zero
+      title: Text-to-Video Zero
+    - local: api/pipelines/unclip
+      title: UnCLIP
+    - local: api/pipelines/latent_diffusion_uncond
+      title: Unconditional Latent Diffusion
+    - local: api/pipelines/versatile_diffusion
+      title: Versatile Diffusion
+    - local: api/pipelines/vq_diffusion
+      title: VQ Diffusion
+    title: Pipelines
+  - sections:
+    - local: api/schedulers/overview
+      title: Overview
+    - local: api/schedulers/ddim
+      title: DDIM
+    - local: api/schedulers/ddim_inverse
+      title: DDIMInverse
+    - local: api/schedulers/ddpm
+      title: DDPM
+    - local: api/schedulers/deis
+      title: DEIS
+    - local: api/schedulers/dpm_discrete
+      title: DPM Discrete Scheduler
+    - local: api/schedulers/dpm_discrete_ancestral
+      title: DPM Discrete Scheduler with ancestral sampling
+    - local: api/schedulers/euler_ancestral
+      title: Euler Ancestral Scheduler
+    - local: api/schedulers/euler
+      title: Euler scheduler
+    - local: api/schedulers/heun
+      title: Heun Scheduler
+    - local: api/schedulers/ipndm
+      title: IPNDM
+    - local: api/schedulers/lms_discrete
+      title: Linear Multistep
+    - local: api/schedulers/multistep_dpm_solver
+      title: Multistep DPM-Solver
+    - local: api/schedulers/pndm
+      title: PNDM
+    - local: api/schedulers/repaint
+      title: RePaint Scheduler
+    - local: api/schedulers/singlestep_dpm_solver
+      title: Singlestep DPM-Solver
+    - local: api/schedulers/stochastic_karras_ve
+      title: Stochastic Kerras VE
+    - local: api/schedulers/unipc
+      title: UniPCMultistepScheduler
+    - local: api/schedulers/score_sde_ve
+      title: VE-SDE
+    - local: api/schedulers/score_sde_vp
+      title: VP-SDE
+    - local: api/schedulers/vq_diffusion
+      title: VQDiffusionScheduler
+    title: Schedulers
+  - sections:
+    - local: api/experimental/rl
+      title: RL Planning
+    title: Experimental Features
+  title: API

examples/community/README.md

@@ -37,7 +37,6 @@ If a community doesn't work as expected, please open an issue and ping the autho
 | TensorRT Stable Diffusion Image to Image Pipeline                                                                                                    | Accelerates the Stable Diffusion Image2Image Pipeline using TensorRT                                                                                                                                                                                                                                                                                                                                                                                                                                      | [TensorRT Stable Diffusion Image to Image Pipeline](#tensorrt-image2image-stable-diffusion-pipeline)      | - |              [Asfiya Baig](https://github.com/asfiyab-nvidia) |
 | Stable Diffusion IPEX Pipeline | Accelerate Stable Diffusion inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [Stable Diffusion on IPEX](#stable-diffusion-on-ipex) | - | [Yingjie Han](https://github.com/yingjie-han/) | 
 | CLIP Guided Images Mixing Stable Diffusion Pipeline | Сombine images using usual diffusion models. | [CLIP Guided Images Mixing Using Stable Diffusion](#clip-guided-images-mixing-with-stable-diffusion) | - | [Karachev Denis](https://github.com/TheDenk) |  
-| TensorRT Stable Diffusion Inpainting Pipeline                                                                                                    | Accelerates the Stable Diffusion Inpainting Pipeline using TensorRT                                                                                                                                                                                                                                                                                                                                                                                                                                      | [TensorRT Stable Diffusion Inpainting Pipeline](#tensorrt-inpainting-stable-diffusion-pipeline)      | - |              [Asfiya Baig](https://github.com/asfiyab-nvidia) |
 
 To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly.
 ```py
@@ -1601,17 +1600,18 @@ pipe_images = mixing_pipeline(
 
 ![image_mixing_result](https://huggingface.co/datasets/TheDenk/images_mixing/resolve/main/boromir_gigachad.png)
 
-### Stable Diffusion Mixture Tiling
+### Stable Diffusion Mixture
 
 This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
     
 ```python
-from diffusers import LMSDiscreteScheduler, DiffusionPipeline
+from diffusers import LMSDiscreteScheduler
+from mixdiff import StableDiffusionTilingPipeline
 
 # Creater scheduler and model (similar to StableDiffusionPipeline)
 scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
-pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler, custom_pipeline="mixture_tiling")
-pipeline.to("cuda")
+pipeline = StableDiffusionTilingPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler)
+pipeline.to("cuda:0")
 
 # Mixture of Diffusers generation
 image = pipeline(
@@ -1629,81 +1629,3 @@ image = pipeline(
     num_inference_steps=50,
 )["images"][0]
 ```
-![mixture_tiling_results](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/mixture_tiling.png)
-
-### TensorRT Inpainting Stable Diffusion Pipeline
-
-The TensorRT Pipeline can be used to accelerate the Inpainting Stable Diffusion Inference run.
-
-NOTE: The ONNX conversions and TensorRT engine build may take up to 30 minutes.
-
-```python
-import requests
-from io import BytesIO
-from PIL import Image
-import torch
-from diffusers import PNDMScheduler
-from diffusers.pipelines.stable_diffusion import StableDiffusionImg2ImgPipeline
-
-# Use the PNDMScheduler scheduler here instead
-scheduler = PNDMScheduler.from_pretrained("stabilityai/stable-diffusion-2-inpainting", subfolder="scheduler")
-
-
-pipe = StableDiffusionImg2ImgPipeline.from_pretrained("stabilityai/stable-diffusion-2-inpainting",
-    custom_pipeline="stable_diffusion_tensorrt_inpaint",
-    revision='fp16',
-    torch_dtype=torch.float16,
-    scheduler=scheduler,
-    )
-
-# re-use cached folder to save ONNX models and TensorRT Engines
-pipe.set_cached_folder("stabilityai/stable-diffusion-2-inpainting", revision='fp16',)
-
-pipe = pipe.to("cuda")
-
-url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
-response = requests.get(url)
-input_image = Image.open(BytesIO(response.content)).convert("RGB")
-
-mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
-response = requests.get(mask_url)
-mask_image = Image.open(BytesIO(response.content)).convert("RGB")
-
-prompt = "a mecha robot sitting on a bench"
-image = pipe(prompt, image=input_image, mask_image=mask_image, strength=0.75,).images[0]
-image.save('tensorrt_inpaint_mecha_robot.png')
-```
-
-### Stable Diffusion Mixture Canvas
-
-This pipeline uses the Mixture. Refer to the [Mixture](https://arxiv.org/abs/2302.02412) paper for more details.
-    
-```python
-from PIL import Image
-from diffusers import LMSDiscreteScheduler, DiffusionPipeline
-from diffusers.pipelines.pipeline_utils import Image2ImageRegion, Text2ImageRegion, preprocess_image
-
-
-# Load and preprocess guide image
-iic_image = preprocess_image(Image.open("input_image.png").convert("RGB"))
-
-# Creater scheduler and model (similar to StableDiffusionPipeline)
-scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
-pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler).to("cuda:0", custom_pipeline="mixture_canvas")
-pipeline.to("cuda")
-
-# Mixture of Diffusers generation
-output = pipeline(
-    canvas_height=800,
-    canvas_width=352,
-    regions=[
-        Text2ImageRegion(0, 800, 0, 352, guidance_scale=8,
-            prompt=f"best quality, masterpiece, WLOP, sakimichan, art contest winner on pixiv, 8K, intricate details, wet effects, rain drops, ethereal, mysterious, futuristic, UHD, HDR, cinematic lighting, in a beautiful forest, rainy day, award winning, trending on artstation, beautiful confident cheerful young woman, wearing a futuristic sleeveless dress, ultra beautiful detailed  eyes, hyper-detailed face, complex,  perfect, model,  textured,  chiaroscuro, professional make-up, realistic, figure in frame, "),
-        Image2ImageRegion(352-800, 352, 0, 352, reference_image=iic_image, strength=1.0),
-    ],
-    num_inference_steps=100,
-    seed=5525475061,
-)["images"][0]
-```
-![Input_Image](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/input_image.png)
-![mixture_canvas_results](https://huggingface.co/datasets/kadirnar/diffusers_readme_images/resolve/main/canvas.png)

examples/community/lpw_stable_diffusion.py

@@ -11,7 +11,7 @@ from transformers import CLIPImageProcessor, CLIPTextModel, CLIPTokenizer
 from diffusers import DiffusionPipeline
 from diffusers.configuration_utils import FrozenDict
 from diffusers.image_processor import VaeImageProcessor
-from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin
+from diffusers.loaders import FromCkptMixin, LoraLoaderMixin, TextualInversionLoaderMixin
 from diffusers.models import AutoencoderKL, UNet2DConditionModel
 from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput, StableDiffusionSafetyChecker
 from diffusers.schedulers import KarrasDiffusionSchedulers
@@ -410,7 +410,7 @@ def preprocess_mask(mask, batch_size, scale_factor=8):
 
 
 class StableDiffusionLongPromptWeightingPipeline(
-    DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromSingleFileMixin
+    DiffusionPipeline, TextualInversionLoaderMixin, LoraLoaderMixin, FromCkptMixin
 ):
     r"""
     Pipeline for text-to-image generation using Stable Diffusion without tokens length limit, and support parsing

examples/community/mixture.py

@@ -4,7 +4,9 @@ from enum import Enum
 from typing import List, Optional, Tuple, Union
 
 import torch
+from ligo.segments import segment
 from tqdm.auto import tqdm
+from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
 
 from diffusers.models import AutoencoderKL, UNet2DConditionModel
 from diffusers.pipeline_utils import DiffusionPipeline
@@ -13,22 +15,17 @@ from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMSchedu
 from diffusers.utils import logging
 
 
-try:
-    from ligo.segments import segment
-    from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
-except ImportError:
-    raise ImportError("Please install transformers and ligo-segments to use the mixture pipeline")
-
 logger = logging.get_logger(__name__)  # pylint: disable=invalid-name
 
 EXAMPLE_DOC_STRING = """
     Examples:
         ```py
-        >>> from diffusers import LMSDiscreteScheduler, DiffusionPipeline
+        >>> from diffusers import LMSDiscreteScheduler
+        >>> from mixdiff import StableDiffusionTilingPipeline
 
         >>> scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
-        >>> pipeline = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler, custom_pipeline="mixture_tiling")
-        >>> pipeline.to("cuda")
+        >>> pipeline = StableDiffusionTilingPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=scheduler)
+        >>> pipeline.to("cuda:0")
 
         >>> image = pipeline(
         >>>     prompt=[[
@@ -218,7 +215,6 @@ class StableDiffusionTilingPipeline(DiffusionPipeline, StableDiffusionExtrasMixi
             raise ValueError(f"`seed_tiles_mode` has to be a string or list of lists but is {type(prompt)}")
         if isinstance(seed_tiles_mode, str):
             seed_tiles_mode = [[seed_tiles_mode for _ in range(len(row))] for row in prompt]
-
         modes = [mode.value for mode in self.SeedTilesMode]
         if any(mode not in modes for row in seed_tiles_mode for mode in row):
             raise ValueError(f"Seed tiles mode must be one of {modes}")

examples/community/mixture_canvas.py

@@ -1,503 +0,0 @@
-import re
-from copy import deepcopy
-from dataclasses import asdict, dataclass
-from enum import Enum
-from typing import List, Optional, Union
-
-import numpy as np
-import torch
-from numpy import exp, pi, sqrt
-from torchvision.transforms.functional import resize
-from tqdm.auto import tqdm
-from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
-
-from diffusers.models import AutoencoderKL, UNet2DConditionModel
-from diffusers.pipeline_utils import DiffusionPipeline
-from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
-from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
-
-
-def preprocess_image(image):
-    from PIL import Image
-
-    """Preprocess an input image
-
-    Same as
-    https://github.com/huggingface/diffusers/blob/1138d63b519e37f0ce04e027b9f4a3261d27c628/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L44
-    """
-    w, h = image.size
-    w, h = (x - x % 32 for x in (w, h))  # resize to integer multiple of 32
-    image = image.resize((w, h), resample=Image.LANCZOS)
-    image = np.array(image).astype(np.float32) / 255.0
-    image = image[None].transpose(0, 3, 1, 2)
-    image = torch.from_numpy(image)
-    return 2.0 * image - 1.0
-
-
-@dataclass
-class CanvasRegion:
-    """Class defining a rectangular region in the canvas"""
-
-    row_init: int  # Region starting row in pixel space (included)
-    row_end: int  # Region end row in pixel space (not included)
-    col_init: int  # Region starting column in pixel space (included)
-    col_end: int  # Region end column in pixel space (not included)
-    region_seed: int = None  # Seed for random operations in this region
-    noise_eps: float = 0.0  # Deviation of a zero-mean gaussian noise to be applied over the latents in this region. Useful for slightly "rerolling" latents
-
-    def __post_init__(self):
-        # Initialize arguments if not specified
-        if self.region_seed is None:
-            self.region_seed = np.random.randint(9999999999)
-        # Check coordinates are non-negative
-        for coord in [self.row_init, self.row_end, self.col_init, self.col_end]:
-            if coord < 0:
-                raise ValueError(
-                    f"A CanvasRegion must be defined with non-negative indices, found ({self.row_init}, {self.row_end}, {self.col_init}, {self.col_end})"
-                )
-        # Check coordinates are divisible by 8, else we end up with nasty rounding error when mapping to latent space
-        for coord in [self.row_init, self.row_end, self.col_init, self.col_end]:
-            if coord // 8 != coord / 8:
-                raise ValueError(
-                    f"A CanvasRegion must be defined with locations divisible by 8, found ({self.row_init}-{self.row_end}, {self.col_init}-{self.col_end})"
-                )
-        # Check noise eps is non-negative
-        if self.noise_eps < 0:
-            raise ValueError(f"A CanvasRegion must be defined noises eps non-negative, found {self.noise_eps}")
-        # Compute coordinates for this region in latent space
-        self.latent_row_init = self.row_init // 8
-        self.latent_row_end = self.row_end // 8
-        self.latent_col_init = self.col_init // 8
-        self.latent_col_end = self.col_end // 8
-
-    @property
-    def width(self):
-        return self.col_end - self.col_init
-
-    @property
-    def height(self):
-        return self.row_end - self.row_init
-
-    def get_region_generator(self, device="cpu"):
-        """Creates a torch.Generator based on the random seed of this region"""
-        # Initialize region generator
-        return torch.Generator(device).manual_seed(self.region_seed)
-
-    @property
-    def __dict__(self):
-        return asdict(self)
-
-
-class MaskModes(Enum):
-    """Modes in which the influence of diffuser is masked"""
-
-    CONSTANT = "constant"
-    GAUSSIAN = "gaussian"
-    QUARTIC = "quartic"  # See https://en.wikipedia.org/wiki/Kernel_(statistics)
-
-
-@dataclass
-class DiffusionRegion(CanvasRegion):
-    """Abstract class defining a region where some class of diffusion process is acting"""
-
-    pass
-
-
-@dataclass
-class Text2ImageRegion(DiffusionRegion):
-    """Class defining a region where a text guided diffusion process is acting"""
-
-    prompt: str = ""  # Text prompt guiding the diffuser in this region
-    guidance_scale: float = 7.5  # Guidance scale of the diffuser in this region. If None, randomize
-    mask_type: MaskModes = MaskModes.GAUSSIAN.value  # Kind of weight mask applied to this region
-    mask_weight: float = 1.0  # Global weights multiplier of the mask
-    tokenized_prompt = None  # Tokenized prompt
-    encoded_prompt = None  # Encoded prompt
-
-    def __post_init__(self):
-        super().__post_init__()
-        # Mask weight cannot be negative
-        if self.mask_weight < 0:
-            raise ValueError(
-                f"A Text2ImageRegion must be defined with non-negative mask weight, found {self.mask_weight}"
-            )
-        # Mask type must be an actual known mask
-        if self.mask_type not in [e.value for e in MaskModes]:
-            raise ValueError(
-                f"A Text2ImageRegion was defined with mask {self.mask_type}, which is not an accepted mask ({[e.value for e in MaskModes]})"
-            )
-        # Randomize arguments if given as None
-        if self.guidance_scale is None:
-            self.guidance_scale = np.random.randint(5, 30)
-        # Clean prompt
-        self.prompt = re.sub(" +", " ", self.prompt).replace("\n", " ")
-
-    def tokenize_prompt(self, tokenizer):
-        """Tokenizes the prompt for this diffusion region using a given tokenizer"""
-        self.tokenized_prompt = tokenizer(
-            self.prompt,
-            padding="max_length",
-            max_length=tokenizer.model_max_length,
-            truncation=True,
-            return_tensors="pt",
-        )
-
-    def encode_prompt(self, text_encoder, device):
-        """Encodes the previously tokenized prompt for this diffusion region using a given encoder"""
-        assert self.tokenized_prompt is not None, ValueError(
-            "Prompt in diffusion region must be tokenized before encoding"
-        )
-        self.encoded_prompt = text_encoder(self.tokenized_prompt.input_ids.to(device))[0]
-
-
-@dataclass
-class Image2ImageRegion(DiffusionRegion):
-    """Class defining a region where an image guided diffusion process is acting"""
-
-    reference_image: torch.FloatTensor = None
-    strength: float = 0.8  # Strength of the image
-
-    def __post_init__(self):
-        super().__post_init__()
-        if self.reference_image is None:
-            raise ValueError("Must provide a reference image when creating an Image2ImageRegion")
-        if self.strength < 0 or self.strength > 1:
-            raise ValueError(f"The value of strength should in [0.0, 1.0] but is {self.strength}")
-        # Rescale image to region shape
-        self.reference_image = resize(self.reference_image, size=[self.height, self.width])
-
-    def encode_reference_image(self, encoder, device, generator, cpu_vae=False):
-        """Encodes the reference image for this Image2Image region into the latent space"""
-        # Place encoder in CPU or not following the parameter cpu_vae
-        if cpu_vae:
-            # Note here we use mean instead of sample, to avoid moving also generator to CPU, which is troublesome
-            self.reference_latents = encoder.cpu().encode(self.reference_image).latent_dist.mean.to(device)
-        else:
-            self.reference_latents = encoder.encode(self.reference_image.to(device)).latent_dist.sample(
-                generator=generator
-            )
-        self.reference_latents = 0.18215 * self.reference_latents
-
-    @property
-    def __dict__(self):
-        # This class requires special casting to dict because of the reference_image tensor. Otherwise it cannot be casted to JSON
-
-        # Get all basic fields from parent class
-        super_fields = {key: getattr(self, key) for key in DiffusionRegion.__dataclass_fields__.keys()}
-        # Pack other fields
-        return {**super_fields, "reference_image": self.reference_image.cpu().tolist(), "strength": self.strength}
-
-
-class RerollModes(Enum):
-    """Modes in which the reroll regions operate"""
-
-    RESET = "reset"  # Completely reset the random noise in the region
-    EPSILON = "epsilon"  # Alter slightly the latents in the region
-
-
-@dataclass
-class RerollRegion(CanvasRegion):
-    """Class defining a rectangular canvas region in which initial latent noise will be rerolled"""
-
-    reroll_mode: RerollModes = RerollModes.RESET.value
-
-
-@dataclass
-class MaskWeightsBuilder:
-    """Auxiliary class to compute a tensor of weights for a given diffusion region"""
-
-    latent_space_dim: int  # Size of the U-net latent space
-    nbatch: int = 1  # Batch size in the U-net
-
-    def compute_mask_weights(self, region: DiffusionRegion) -> torch.tensor:
-        """Computes a tensor of weights for a given diffusion region"""
-        MASK_BUILDERS = {
-            MaskModes.CONSTANT.value: self._constant_weights,
-            MaskModes.GAUSSIAN.value: self._gaussian_weights,
-            MaskModes.QUARTIC.value: self._quartic_weights,
-        }
-        return MASK_BUILDERS[region.mask_type](region)
-
-    def _constant_weights(self, region: DiffusionRegion) -> torch.tensor:
-        """Computes a tensor of constant for a given diffusion region"""
-        latent_width = region.latent_col_end - region.latent_col_init
-        latent_height = region.latent_row_end - region.latent_row_init
-        return torch.ones(self.nbatch, self.latent_space_dim, latent_height, latent_width) * region.mask_weight
-
-    def _gaussian_weights(self, region: DiffusionRegion) -> torch.tensor:
-        """Generates a gaussian mask of weights for tile contributions"""
-        latent_width = region.latent_col_end - region.latent_col_init
-        latent_height = region.latent_row_end - region.latent_row_init
-
-        var = 0.01
-        midpoint = (latent_width - 1) / 2  # -1 because index goes from 0 to latent_width - 1
-        x_probs = [
-            exp(-(x - midpoint) * (x - midpoint) / (latent_width * latent_width) / (2 * var)) / sqrt(2 * pi * var)
-            for x in range(latent_width)
-        ]
-        midpoint = (latent_height - 1) / 2
-        y_probs = [
-            exp(-(y - midpoint) * (y - midpoint) / (latent_height * latent_height) / (2 * var)) / sqrt(2 * pi * var)
-            for y in range(latent_height)
-        ]
-
-        weights = np.outer(y_probs, x_probs) * region.mask_weight
-        return torch.tile(torch.tensor(weights), (self.nbatch, self.latent_space_dim, 1, 1))
-
-    def _quartic_weights(self, region: DiffusionRegion) -> torch.tensor:
-        """Generates a quartic mask of weights for tile contributions
-
-        The quartic kernel has bounded support over the diffusion region, and a smooth decay to the region limits.
-        """
-        quartic_constant = 15.0 / 16.0
-
-        support = (np.array(range(region.latent_col_init, region.latent_col_end)) - region.latent_col_init) / (
-            region.latent_col_end - region.latent_col_init - 1
-        ) * 1.99 - (1.99 / 2.0)
-        x_probs = quartic_constant * np.square(1 - np.square(support))
-        support = (np.array(range(region.latent_row_init, region.latent_row_end)) - region.latent_row_init) / (
-            region.latent_row_end - region.latent_row_init - 1
-        ) * 1.99 - (1.99 / 2.0)
-        y_probs = quartic_constant * np.square(1 - np.square(support))
-
-        weights = np.outer(y_probs, x_probs) * region.mask_weight
-        return torch.tile(torch.tensor(weights), (self.nbatch, self.latent_space_dim, 1, 1))
-
-
-class StableDiffusionCanvasPipeline(DiffusionPipeline):
-    """Stable Diffusion pipeline that mixes several diffusers in the same canvas"""
-
-    def __init__(
-        self,
-        vae: AutoencoderKL,
-        text_encoder: CLIPTextModel,
-        tokenizer: CLIPTokenizer,
-        unet: UNet2DConditionModel,
-        scheduler: Union[DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler],
-        safety_checker: StableDiffusionSafetyChecker,
-        feature_extractor: CLIPFeatureExtractor,
-    ):
-        super().__init__()
-        self.register_modules(
-            vae=vae,
-            text_encoder=text_encoder,
-            tokenizer=tokenizer,
-            unet=unet,
-            scheduler=scheduler,
-            safety_checker=safety_checker,
-            feature_extractor=feature_extractor,
-        )
-
-    def decode_latents(self, latents, cpu_vae=False):
-        """Decodes a given array of latents into pixel space"""
-        # scale and decode the image latents with vae
-        if cpu_vae:
-            lat = deepcopy(latents).cpu()
-            vae = deepcopy(self.vae).cpu()
-        else:
-            lat = latents
-            vae = self.vae
-
-        lat = 1 / 0.18215 * lat
-        image = vae.decode(lat).sample
-
-        image = (image / 2 + 0.5).clamp(0, 1)
-        image = image.cpu().permute(0, 2, 3, 1).numpy()
-
-        return self.numpy_to_pil(image)
-
-    def get_latest_timestep_img2img(self, num_inference_steps, strength):
-        """Finds the latest timesteps where an img2img strength does not impose latents anymore"""
-        # get the original timestep using init_timestep
-        offset = self.scheduler.config.get("steps_offset", 0)
-        init_timestep = int(num_inference_steps * (1 - strength)) + offset
-        init_timestep = min(init_timestep, num_inference_steps)
-
-        t_start = min(max(num_inference_steps - init_timestep + offset, 0), num_inference_steps - 1)
-        latest_timestep = self.scheduler.timesteps[t_start]
-
-        return latest_timestep
-
-    @torch.no_grad()
-    def __call__(
-        self,
-        canvas_height: int,
-        canvas_width: int,
-        regions: List[DiffusionRegion],
-        num_inference_steps: Optional[int] = 50,
-        seed: Optional[int] = 12345,
-        reroll_regions: Optional[List[RerollRegion]] = None,
-        cpu_vae: Optional[bool] = False,
-        decode_steps: Optional[bool] = False,
-    ):
-        if reroll_regions is None:
-            reroll_regions = []
-        batch_size = 1
-
-        if decode_steps:
-            steps_images = []
-
-        # Prepare scheduler
-        self.scheduler.set_timesteps(num_inference_steps, device=self.device)
-
-        # Split diffusion regions by their kind
-        text2image_regions = [region for region in regions if isinstance(region, Text2ImageRegion)]
-        image2image_regions = [region for region in regions if isinstance(region, Image2ImageRegion)]
-
-        # Prepare text embeddings
-        for region in text2image_regions:
-            region.tokenize_prompt(self.tokenizer)
-            region.encode_prompt(self.text_encoder, self.device)
-
-        # Create original noisy latents using the timesteps
-        latents_shape = (batch_size, self.unet.config.in_channels, canvas_height // 8, canvas_width // 8)
-        generator = torch.Generator(self.device).manual_seed(seed)
-        init_noise = torch.randn(latents_shape, generator=generator, device=self.device)
-
-        # Reset latents in seed reroll regions, if requested
-        for region in reroll_regions:
-            if region.reroll_mode == RerollModes.RESET.value:
-                region_shape = (
-                    latents_shape[0],
-                    latents_shape[1],
-                    region.latent_row_end - region.latent_row_init,
-                    region.latent_col_end - region.latent_col_init,
-                )
-                init_noise[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ] = torch.randn(region_shape, generator=region.get_region_generator(self.device), device=self.device)
-
-        # Apply epsilon noise to regions: first diffusion regions, then reroll regions
-        all_eps_rerolls = regions + [r for r in reroll_regions if r.reroll_mode == RerollModes.EPSILON.value]
-        for region in all_eps_rerolls:
-            if region.noise_eps > 0:
-                region_noise = init_noise[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ]
-                eps_noise = (
-                    torch.randn(
-                        region_noise.shape, generator=region.get_region_generator(self.device), device=self.device
-                    )
-                    * region.noise_eps
-                )
-                init_noise[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ] += eps_noise
-
-        # scale the initial noise by the standard deviation required by the scheduler
-        latents = init_noise * self.scheduler.init_noise_sigma
-
-        # Get unconditional embeddings for classifier free guidance in text2image regions
-        for region in text2image_regions:
-            max_length = region.tokenized_prompt.input_ids.shape[-1]
-            uncond_input = self.tokenizer(
-                [""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt"
-            )
-            uncond_embeddings = self.text_encoder(uncond_input.input_ids.to(self.device))[0]
-
-            # For classifier free guidance, we need to do two forward passes.
-            # Here we concatenate the unconditional and text embeddings into a single batch
-            # to avoid doing two forward passes
-            region.encoded_prompt = torch.cat([uncond_embeddings, region.encoded_prompt])
-
-        # Prepare image latents
-        for region in image2image_regions:
-            region.encode_reference_image(self.vae, device=self.device, generator=generator)
-
-        # Prepare mask of weights for each region
-        mask_builder = MaskWeightsBuilder(latent_space_dim=self.unet.config.in_channels, nbatch=batch_size)
-        mask_weights = [mask_builder.compute_mask_weights(region).to(self.device) for region in text2image_regions]
-
-        # Diffusion timesteps
-        for i, t in tqdm(enumerate(self.scheduler.timesteps)):
-            # Diffuse each region
-            noise_preds_regions = []
-
-            # text2image regions
-            for region in text2image_regions:
-                region_latents = latents[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ]
-                # expand the latents if we are doing classifier free guidance
-                latent_model_input = torch.cat([region_latents] * 2)
-                # scale model input following scheduler rules
-                latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
-                # predict the noise residual
-                noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=region.encoded_prompt)["sample"]
-                # perform guidance
-                noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
-                noise_pred_region = noise_pred_uncond + region.guidance_scale * (noise_pred_text - noise_pred_uncond)
-                noise_preds_regions.append(noise_pred_region)
-
-            # Merge noise predictions for all tiles
-            noise_pred = torch.zeros(latents.shape, device=self.device)
-            contributors = torch.zeros(latents.shape, device=self.device)
-            # Add each tile contribution to overall latents
-            for region, noise_pred_region, mask_weights_region in zip(
-                text2image_regions, noise_preds_regions, mask_weights
-            ):
-                noise_pred[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ] += (
-                    noise_pred_region * mask_weights_region
-                )
-                contributors[
-                    :,
-                    :,
-                    region.latent_row_init : region.latent_row_end,
-                    region.latent_col_init : region.latent_col_end,
-                ] += mask_weights_region
-            # Average overlapping areas with more than 1 contributor
-            noise_pred /= contributors
-            noise_pred = torch.nan_to_num(
-                noise_pred
-            )  # Replace NaNs by zeros: NaN can appear if a position is not covered by any DiffusionRegion
-
-            # compute the previous noisy sample x_t -> x_t-1
-            latents = self.scheduler.step(noise_pred, t, latents).prev_sample
-
-            # Image2Image regions: override latents generated by the scheduler
-            for region in image2image_regions:
-                influence_step = self.get_latest_timestep_img2img(num_inference_steps, region.strength)
-                # Only override in the timesteps before the last influence step of the image (given by its strength)
-                if t > influence_step:
-                    timestep = t.repeat(batch_size)
-                    region_init_noise = init_noise[
-                        :,
-                        :,
-                        region.latent_row_init : region.latent_row_end,
-                        region.latent_col_init : region.latent_col_end,
-                    ]
-                    region_latents = self.scheduler.add_noise(region.reference_latents, region_init_noise, timestep)
-                    latents[
-                        :,
-                        :,
-                        region.latent_row_init : region.latent_row_end,
-                        region.latent_col_init : region.latent_col_end,
-                    ] = region_latents
-
-            if decode_steps:
-                steps_images.append(self.decode_latents(latents, cpu_vae))
-
-        # scale and decode the image latents with vae
-        image = self.decode_latents(latents, cpu_vae)
-
-        output = {"images": image}
-        if decode_steps:
-            output = {**output, "steps_images": steps_images}
-        return output

examples/community/stable_diffusion_controlnet_reference.py

@@ -1,7 +1,6 @@
 # Inspired by: https://github.com/Mikubill/sd-webui-controlnet/discussions/1236 and https://github.com/Mikubill/sd-webui-controlnet/discussions/1280
 from typing import Any, Callable, Dict, List, Optional, Tuple, Union
 
-import numpy as np
 import PIL.Image
 import torch
 
@@ -98,14 +97,7 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
     def __call__(
         self,
         prompt: Union[str, List[str]] = None,
-        image: Union[
-            torch.FloatTensor,
-            PIL.Image.Image,
-            np.ndarray,
-            List[torch.FloatTensor],
-            List[PIL.Image.Image],
-            List[np.ndarray],
-        ] = None,
+        image: Union[torch.FloatTensor, PIL.Image.Image, List[torch.FloatTensor], List[PIL.Image.Image]] = None,
         ref_image: Union[torch.FloatTensor, PIL.Image.Image] = None,
         height: Optional[int] = None,
         width: Optional[int] = None,
@@ -138,8 +130,8 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
             prompt (`str` or `List[str]`, *optional*):
                 The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`.
                 instead.
-            image (`torch.FloatTensor`, `PIL.Image.Image`, `np.ndarray`, `List[torch.FloatTensor]`, `List[PIL.Image.Image]`, `List[np.ndarray]`,:
-                    `List[List[torch.FloatTensor]]`, `List[List[np.ndarray]]` or `List[List[PIL.Image.Image]]`):
+            image (`torch.FloatTensor`, `PIL.Image.Image`, `List[torch.FloatTensor]`, `List[PIL.Image.Image]`,
+                    `List[List[torch.FloatTensor]]`, or `List[List[PIL.Image.Image]]`):
                 The ControlNet input condition. ControlNet uses this input condition to generate guidance to Unet. If
                 the type is specified as `Torch.FloatTensor`, it is passed to ControlNet as is. `PIL.Image.Image` can
                 also be accepted as an image. The dimensions of the output image defaults to `image`'s dimensions. If
@@ -231,12 +223,15 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
             list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
             (nsfw) content, according to the `safety_checker`.
         """
-        assert reference_attn or reference_adain, "`reference_attn` or `reference_adain` must be True."
+        # 0. Default height and width to unet
+        height, width = self._default_height_width(height, width, image)
 
         # 1. Check inputs. Raise error if not correct
         self.check_inputs(
             prompt,
             image,
+            height,
+            width,
             callback_steps,
             negative_prompt,
             prompt_embeds,
@@ -271,9 +266,6 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
         guess_mode = guess_mode or global_pool_conditions
 
         # 3. Encode input prompt
-        text_encoder_lora_scale = (
-            cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None
-        )
         prompt_embeds = self._encode_prompt(
             prompt,
             device,
@@ -282,7 +274,6 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
             negative_prompt,
             prompt_embeds=prompt_embeds,
             negative_prompt_embeds=negative_prompt_embeds,
-            lora_scale=text_encoder_lora_scale,
         )
 
         # 4. Prepare image
@@ -298,7 +289,6 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
                 do_classifier_free_guidance=do_classifier_free_guidance,
                 guess_mode=guess_mode,
             )
-            height, width = image.shape[-2:]
         elif isinstance(controlnet, MultiControlNetModel):
             images = []
 
@@ -318,7 +308,6 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
                 images.append(image_)
 
             image = images
-            height, width = image[0].shape[-2:]
         else:
             assert False
 
@@ -731,15 +720,14 @@ class StableDiffusionControlNetReferencePipeline(StableDiffusionControlNetPipeli
                 # controlnet(s) inference
                 if guess_mode and do_classifier_free_guidance:
                     # Infer ControlNet only for the conditional batch.
-                    control_model_input = latents
-                    control_model_input = self.scheduler.scale_model_input(control_model_input, t)
+                    controlnet_latent_model_input = latents
                     controlnet_prompt_embeds = prompt_embeds.chunk(2)[1]
                 else:
-                    control_model_input = latent_model_input
+                    controlnet_latent_model_input = latent_model_input
                     controlnet_prompt_embeds = prompt_embeds
 
                 down_block_res_samples, mid_block_res_sample = self.controlnet(
-                    control_model_input,
+                    controlnet_latent_model_input,
                     t,
                     encoder_hidden_states=controlnet_prompt_embeds,
                     controlnet_cond=image,

examples/community/stable_diffusion_reference.py

@@ -9,7 +9,6 @@ from diffusers import StableDiffusionPipeline
 from diffusers.models.attention import BasicTransformerBlock
 from diffusers.models.unet_2d_blocks import CrossAttnDownBlock2D, CrossAttnUpBlock2D, DownBlock2D, UpBlock2D
 from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
-from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion import rescale_noise_cfg
 from diffusers.utils import PIL_INTERPOLATION, logging, randn_tensor
 
 
@@ -180,7 +179,6 @@ class StableDiffusionReferencePipeline(StableDiffusionPipeline):
         callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
         callback_steps: int = 1,
         cross_attention_kwargs: Optional[Dict[str, Any]] = None,
-        guidance_rescale: float = 0.0,
         attention_auto_machine_weight: float = 1.0,
         gn_auto_machine_weight: float = 1.0,
         style_fidelity: float = 0.5,
@@ -250,11 +248,6 @@ class StableDiffusionReferencePipeline(StableDiffusionPipeline):
                 A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under
                 `self.processor` in
                 [diffusers.cross_attention](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py).
-            guidance_rescale (`float`, *optional*, defaults to 0.7):
-                Guidance rescale factor proposed by [Common Diffusion Noise Schedules and Sample Steps are
-                Flawed](https://arxiv.org/pdf/2305.08891.pdf) `guidance_scale` is defined as `φ` in equation 16. of
-                [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf).
-                Guidance rescale factor should fix overexposure when using zero terminal SNR.
             attention_auto_machine_weight (`float`):
                 Weight of using reference query for self attention's context.
                 If attention_auto_machine_weight=1.0, use reference query for all self attention's context.
@@ -302,9 +295,6 @@ class StableDiffusionReferencePipeline(StableDiffusionPipeline):
         do_classifier_free_guidance = guidance_scale > 1.0
 
         # 3. Encode input prompt
-        text_encoder_lora_scale = (
-            cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None
-        )
         prompt_embeds = self._encode_prompt(
             prompt,
             device,
@@ -313,7 +303,6 @@ class StableDiffusionReferencePipeline(StableDiffusionPipeline):
             negative_prompt,
             prompt_embeds=prompt_embeds,
             negative_prompt_embeds=negative_prompt_embeds,
-            lora_scale=text_encoder_lora_scale,
         )
 
         # 4. Preprocess reference image
@@ -759,10 +748,6 @@ class StableDiffusionReferencePipeline(StableDiffusionPipeline):
                     noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
                     noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
 
-                if do_classifier_free_guidance and guidance_rescale > 0.0:
-                    # Based on 3.4. in https://arxiv.org/pdf/2305.08891.pdf
-                    noise_pred = rescale_noise_cfg(noise_pred, noise_pred_text, guidance_rescale=guidance_rescale)
-
                 # compute the previous noisy sample x_t -> x_t-1
                 latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0]
 

examples/community/unclip_image_interpolation.py

@@ -376,16 +376,14 @@ class UnCLIPImageInterpolationPipeline(DiffusionPipeline):
         height = self.decoder.config.sample_size
         width = self.decoder.config.sample_size
 
-        # Get the decoder latents for 1 step and then repeat the same tensor for the entire batch to keep same noise across all interpolation steps.
         decoder_latents = self.prepare_latents(
-            (1, num_channels_latents, height, width),
+            (batch_size, num_channels_latents, height, width),
             text_encoder_hidden_states.dtype,
             device,
             generator,
             decoder_latents,
             self.decoder_scheduler,
         )
-        decoder_latents = decoder_latents.repeat((batch_size, 1, 1, 1))
 
         for i, t in enumerate(self.progress_bar(decoder_timesteps_tensor)):
             # expand the latents if we are doing classifier free guidance

examples/controlnet/train_controlnet.py

@@ -18,7 +18,6 @@ import logging
 import math
 import os
 import random
-import shutil
 from pathlib import Path
 
 import accelerate
@@ -56,7 +55,7 @@ if is_wandb_available():
     import wandb
 
 # Will error if the minimal version of diffusers is not installed. Remove at your own risks.
-check_min_version("0.18.0")
+check_min_version("0.17.0.dev0")
 
 logger = get_logger(__name__)
 
@@ -308,7 +307,11 @@ def parse_args(input_args=None):
         "--checkpoints_total_limit",
         type=int,
         default=None,
-        help=("Max number of checkpoints to store."),
+        help=(
+            "Max number of checkpoints to store. Passed as `total_limit` to the `Accelerator` `ProjectConfiguration`."
+            " See Accelerator::save_state https://huggingface.co/docs/accelerate/package_reference/accelerator#accelerate.Accelerator.save_state"
+            " for more details"
+        ),
     )
     parser.add_argument(
         "--resume_from_checkpoint",
@@ -713,12 +716,13 @@ def collate_fn(examples):
 def main(args):
     logging_dir = Path(args.output_dir, args.logging_dir)
 
-    accelerator_project_config = ProjectConfiguration(project_dir=args.output_dir, logging_dir=logging_dir)
+    accelerator_project_config = ProjectConfiguration(total_limit=args.checkpoints_total_limit)
 
     accelerator = Accelerator(
         gradient_accumulation_steps=args.gradient_accumulation_steps,
         mixed_precision=args.mixed_precision,
         log_with=args.report_to,
+        logging_dir=logging_dir,
         project_config=accelerator_project_config,
     )
 
@@ -1055,26 +1059,6 @@ def main(args):
 
                 if accelerator.is_main_process:
                     if global_step % args.checkpointing_steps == 0:
-                        # _before_ saving state, check if this save would set us over the `checkpoints_total_limit`
-                        if args.checkpoints_total_limit is not None:
-                            checkpoints = os.listdir(args.output_dir)
-                            checkpoints = [d for d in checkpoints if d.startswith("checkpoint")]
-                            checkpoints = sorted(checkpoints, key=lambda x: int(x.split("-")[1]))
-
-                            # before we save the new checkpoint, we need to have at _most_ `checkpoints_total_limit - 1` checkpoints
-                            if len(checkpoints) >= args.checkpoints_total_limit:
-                                num_to_remove = len(checkpoints) - args.checkpoints_total_limit + 1
-                                removing_checkpoints = checkpoints[0:num_to_remove]
-
-                                logger.info(
-                                    f"{len(checkpoints)} checkpoints already exist, removing {len(removing_checkpoints)} checkpoints"
-                                )
-                                logger.info(f"removing checkpoints: {', '.join(removing_checkpoints)}")
-
-                                for removing_checkpoint in removing_checkpoints:
-                                    removing_checkpoint = os.path.join(args.output_dir, removing_checkpoint)
-                                    shutil.rmtree(removing_checkpoint)
-
                         save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
                         accelerator.save_state(save_path)
                         logger.info(f"Saved state to {save_path}")

examples/controlnet/train_controlnet_flax.py

@@ -59,7 +59,7 @@ if is_wandb_available():
     import wandb
 
 # Will error if the minimal version of diffusers is not installed. Remove at your own risks.
-check_min_version("0.18.0")
+check_min_version("0.17.0.dev0")
 
 logger = logging.getLogger(__name__)
 

examples/custom_diffusion/train_custom_diffusion.py

@@ -21,7 +21,6 @@ import logging
 import math
 import os
 import random
-import shutil
 import warnings
 from pathlib import Path
 
@@ -57,7 +56,7 @@ from diffusers.utils.import_utils import is_xformers_available
 
 
 # Will error if the minimal version of diffusers is not installed. Remove at your own risks.
-check_min_version("0.18.0")
+check_min_version("0.17.0.dev0")
 
 logger = get_logger(__name__)
 
@@ -447,7 +446,11 @@ def parse_args(input_args=None):
         "--checkpoints_total_limit",
         type=int,
         default=None,
-        help=("Max number of checkpoints to store."),
+        help=(
+            "Max number of checkpoints to store. Passed as `total_limit` to the `Accelerator` `ProjectConfiguration`."
+            " See Accelerator::save_state https://huggingface.co/docs/accelerate/package_reference/accelerator#accelerate.Accelerator.save_state"
+            " for more docs"
+        ),
     )
     parser.add_argument(
         "--resume_from_checkpoint",
@@ -634,12 +637,13 @@ def parse_args(input_args=None):
 def main(args):
     logging_dir = Path(args.output_dir, args.logging_dir)
 
-    accelerator_project_config = ProjectConfiguration(project_dir=args.output_dir, logging_dir=logging_dir)
+    accelerator_project_config = ProjectConfiguration(total_limit=args.checkpoints_total_limit)
 
     accelerator = Accelerator(
         gradient_accumulation_steps=args.gradient_accumulation_steps,
         mixed_precision=args.mixed_precision,
         log_with=args.report_to,
+        logging_dir=logging_dir,
         project_config=accelerator_project_config,
     )
 
@@ -1166,26 +1170,6 @@ def main(args):
 
                 if global_step % args.checkpointing_steps == 0:
                     if accelerator.is_main_process:
-                        # _before_ saving state, check if this save would set us over the `checkpoints_total_limit`
-                        if args.checkpoints_total_limit is not None:
-                            checkpoints = os.listdir(args.output_dir)
-                            checkpoints = [d for d in checkpoints if d.startswith("checkpoint")]
-                            checkpoints = sorted(checkpoints, key=lambda x: int(x.split("-")[1]))
-
-                            # before we save the new checkpoint, we need to have at _most_ `checkpoints_total_limit - 1` checkpoints
-                            if len(checkpoints) >= args.checkpoints_total_limit:
-                                num_to_remove = len(checkpoints) - args.checkpoints_total_limit + 1
-                                removing_checkpoints = checkpoints[0:num_to_remove]
-
-                                logger.info(
-                                    f"{len(checkpoints)} checkpoints already exist, removing {len(removing_checkpoints)} checkpoints"
-                                )
-                                logger.info(f"removing checkpoints: {', '.join(removing_checkpoints)}")
-
-                                for removing_checkpoint in removing_checkpoints:
-                                    removing_checkpoint = os.path.join(args.output_dir, removing_checkpoint)
-                                    shutil.rmtree(removing_checkpoint)
-
                         save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
                         accelerator.save_state(save_path)
                         logger.info(f"Saved state to {save_path}")

examples/dreambooth/README.md

@@ -536,68 +536,9 @@ You can refer to [this blog post](https://huggingface.co/blog/dreambooth) that d
 
 ## IF
 
-You can use the lora and full dreambooth scripts to train the text to image [IF model](https://huggingface.co/DeepFloyd/IF-I-XL-v1.0) and the stage II upscaler 
-[IF model](https://huggingface.co/DeepFloyd/IF-II-L-v1.0).
+You can use the lora and full dreambooth scripts to also train the text to image [IF model](https://huggingface.co/DeepFloyd/IF-I-XL-v1.0). A few alternative cli flags are needed due to the model size, the expected input resolution, and the text encoder conventions.
 
-Note that IF has a predicted variance, and our finetuning scripts only train the models predicted error, so for finetuned IF models we switch to a fixed
-variance schedule. The full finetuning scripts will update the scheduler config for the full saved model. However, when loading saved LoRA weights, you
-must also update the pipeline's scheduler config.
-
-```py
-from diffusers import DiffusionPipeline
-
-pipe = DiffusionPipeline.from_pretrained("DeepFloyd/IF-I-XL-v1.0")
-
-pipe.load_lora_weights("<lora weights path>")
-
-# Update scheduler config to fixed variance schedule
-pipe.scheduler = pipe.scheduler.__class__.from_config(pipe.scheduler.config, variance_type="fixed_small")
-```
-
-Additionally, a few alternative cli flags are needed for IF.
-
-`--resolution=64`: IF is a pixel space diffusion model. In order to operate on un-compressed pixels, the input images are of a much smaller resolution. 
-
-`--pre_compute_text_embeddings`: IF uses [T5](https://huggingface.co/docs/transformers/model_doc/t5) for its text encoder. In order to save GPU memory, we pre compute all text embeddings and then de-allocate
-T5.
-
-`--tokenizer_max_length=77`: T5 has a longer default text length, but the default IF encoding procedure uses a smaller number.
-
-`--text_encoder_use_attention_mask`: T5 passes the attention mask to the text encoder.
-
-### Tips and Tricks
-We find LoRA to be sufficient for finetuning the stage I model as the low resolution of the model makes representing finegrained detail hard regardless.
-
-For common and/or not-visually complex object concepts, you can get away with not-finetuning the upscaler. Just be sure to adjust the prompt passed to the
-upscaler to remove the new token from the instance prompt. I.e. if your stage I prompt is "a sks dog", use "a dog" for your stage II prompt.
-
-For finegrained detail like faces that aren't present in the original training set, we find that full finetuning of the stage II upscaler is better than 
-LoRA finetuning stage II.
-
-For finegrained detail like faces, we find that lower learning rates along with larger batch sizes work best.
-
-For stage II, we find that lower learning rates are also needed.
-
-We found experimentally that the DDPM scheduler with the default larger number of denoising steps to sometimes work better than the DPM Solver scheduler
-used in the training scripts.
-
-### Stage II additional validation images
-
-The stage II validation requires images to upscale, we can download a downsized version of the training set:
-
-```py
-from huggingface_hub import snapshot_download
-
-local_dir = "./dog_downsized"
-snapshot_download(
-    "diffusers/dog-example-downsized",
-    local_dir=local_dir,
-    repo_type="dataset",
-    ignore_patterns=".gitattributes",
-)
-```
-
-### IF stage I LoRA Dreambooth
+### LoRA Dreambooth
 This training configuration requires ~28 GB VRAM.
 
 ```sh
@@ -611,7 +552,7 @@ accelerate launch train_dreambooth_lora.py \
   --instance_data_dir=$INSTANCE_DIR \
   --output_dir=$OUTPUT_DIR \
   --instance_prompt="a sks dog" \
-  --resolution=64 \
+  --resolution=64 \ # The input resolution of the IF unet is 64x64
   --train_batch_size=4 \
   --gradient_accumulation_steps=1 \
   --learning_rate=5e-6 \
@@ -620,58 +561,16 @@ accelerate launch train_dreambooth_lora.py \
   --validation_prompt="a sks dog" \
   --validation_epochs=25 \
   --checkpointing_steps=100 \
-  --pre_compute_text_embeddings \
-  --tokenizer_max_length=77 \
-  --text_encoder_use_attention_mask
+  --pre_compute_text_embeddings \ # Pre compute text embeddings to that T5 doesn't have to be kept in memory
+  --tokenizer_max_length=77 \ # IF expects an override of the max token length
+  --text_encoder_use_attention_mask # IF expects attention mask for text embeddings
 ```
 
-### IF stage II LoRA Dreambooth
+### Full Dreambooth
+Due to the size of the optimizer states, we recommend training the full XL IF model with 8bit adam. 
+Using 8bit adam and the rest of the following config, the model can be trained in ~48 GB VRAM.
 
-`--validation_images`: These images are upscaled during validation steps.
-
-`--class_labels_conditioning=timesteps`: Pass additional conditioning to the UNet needed for stage II.
-
-`--learning_rate=1e-6`: Lower learning rate than stage I.
-
-`--resolution=256`: The upscaler expects higher resolution inputs
-
-```sh
-export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
-export INSTANCE_DIR="dog"
-export OUTPUT_DIR="dreambooth_dog_upscale"
-export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
-
-python train_dreambooth_lora.py \
-    --report_to wandb \
-    --pretrained_model_name_or_path=$MODEL_NAME \
-    --instance_data_dir=$INSTANCE_DIR \
-    --output_dir=$OUTPUT_DIR \
-    --instance_prompt="a sks dog" \
-    --resolution=256 \
-    --train_batch_size=4 \
-    --gradient_accumulation_steps=1 \
-    --learning_rate=1e-6 \ 
-    --max_train_steps=2000 \
-    --validation_prompt="a sks dog" \
-    --validation_epochs=100 \
-    --checkpointing_steps=500 \
-    --pre_compute_text_embeddings \
-    --tokenizer_max_length=77 \
-    --text_encoder_use_attention_mask \
-    --validation_images $VALIDATION_IMAGES \
-    --class_labels_conditioning=timesteps
-```
-
-### IF Stage I Full Dreambooth
-`--skip_save_text_encoder`: When training the full model, this will skip saving the entire T5 with the finetuned model. You can still load the pipeline
-with a T5 loaded from the original model.
-
-`use_8bit_adam`: Due to the size of the optimizer states, we recommend training the full XL IF model with 8bit adam. 
-
-`--learning_rate=1e-7`: For full dreambooth, IF requires very low learning rates. With higher learning rates model quality will degrade. Note that it is 
-likely the learning rate can be increased with larger batch sizes.
-
-Using 8bit adam and a batch size of 4, the model can be trained in ~48 GB VRAM.
+For full dreambooth, IF requires very low learning rates. With higher learning rates model quality will degrade.
 
 ```sh
 export MODEL_NAME="DeepFloyd/IF-I-XL-v1.0"
@@ -684,56 +583,17 @@ accelerate launch train_dreambooth.py \
   --instance_data_dir=$INSTANCE_DIR \
   --output_dir=$OUTPUT_DIR \
   --instance_prompt="a photo of sks dog" \
-  --resolution=64 \
+  --resolution=64 \ # The input resolution of the IF unet is 64x64
   --train_batch_size=4 \
   --gradient_accumulation_steps=1 \
   --learning_rate=1e-7 \
   --max_train_steps=150 \
   --validation_prompt "a photo of sks dog" \
   --validation_steps 25 \
-  --text_encoder_use_attention_mask \
-  --tokenizer_max_length 77 \
-  --pre_compute_text_embeddings \
-  --use_8bit_adam \
+  --text_encoder_use_attention_mask \ # IF expects attention mask for text embeddings
+  --tokenizer_max_length 77 \ # IF expects an override of the max token length
+  --pre_compute_text_embeddings \ # Pre compute text embeddings to that T5 doesn't have to be kept in memory
+  --use_8bit_adam \ # 
   --set_grads_to_none \
-  --skip_save_text_encoder \
-  --push_to_hub
-```
-
-### IF Stage II Full Dreambooth
-
-`--learning_rate=5e-6`: With a smaller effective batch size of 4, we found that we required learning rates as low as
-1e-8.
-
-`--resolution=256`: The upscaler expects higher resolution inputs
-
-`--train_batch_size=2` and `--gradient_accumulation_steps=6`: We found that full training of stage II particularly with
-faces required large effective batch sizes.
-
-```sh
-export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
-export INSTANCE_DIR="dog"
-export OUTPUT_DIR="dreambooth_dog_upscale"
-export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
-
-accelerate launch train_dreambooth.py \
-  --report_to wandb \
-  --pretrained_model_name_or_path=$MODEL_NAME \
-  --instance_data_dir=$INSTANCE_DIR \
-  --output_dir=$OUTPUT_DIR \
-  --instance_prompt="a sks dog" \
-  --resolution=256 \
-  --train_batch_size=2 \
-  --gradient_accumulation_steps=6 \
-  --learning_rate=5e-6 \
-  --max_train_steps=2000 \
-  --validation_prompt="a sks dog" \
-  --validation_steps=150 \
-  --checkpointing_steps=500 \
-  --pre_compute_text_embeddings \
-  --tokenizer_max_length=77 \
-  --text_encoder_use_attention_mask \
-  --validation_images $VALIDATION_IMAGES \
-  --class_labels_conditioning timesteps \
-  --push_to_hub
+  --skip_save_text_encoder # do not save the full T5 text encoder with the model
 ```

examples/dreambooth/train_dreambooth.py

@@ -20,7 +20,6 @@ import itertools
 import logging
 import math
 import os
-import shutil
 import warnings
 from pathlib import Path
 
@@ -59,7 +58,7 @@ if is_wandb_available():
     import wandb
 
 # Will error if the minimal version of diffusers is not installed. Remove at your own risks.
-check_min_version("0.18.0")
+check_min_version("0.17.0.dev0")
 
 logger = get_logger(__name__)
 
@@ -115,17 +114,16 @@ def log_validation(
 
     pipeline_args = {}
 
+    if text_encoder is not None:
+        pipeline_args["text_encoder"] = accelerator.unwrap_model(text_encoder)
+
     if vae is not None:
         pipeline_args["vae"] = vae
 
-    if text_encoder is not None:
-        text_encoder = accelerator.unwrap_model(text_encoder)
-
     # create pipeline (note: unet and vae are loaded again in float32)
     pipeline = DiffusionPipeline.from_pretrained(
         args.pretrained_model_name_or_path,
         tokenizer=tokenizer,
-        text_encoder=text_encoder,
         unet=accelerator.unwrap_model(unet),
         revision=args.revision,
         torch_dtype=weight_dtype,
@@ -158,16 +156,10 @@ def log_validation(
     # run inference
     generator = None if args.seed is None else torch.Generator(device=accelerator.device).manual_seed(args.seed)
     images = []
-    if args.validation_images is None:
-        for _ in range(args.num_validation_images):
-            with torch.autocast("cuda"):
-                image = pipeline(**pipeline_args, num_inference_steps=25, generator=generator).images[0]
-            images.append(image)
-    else:
-        for image in args.validation_images:
-            image = Image.open(image)
-            image = pipeline(**pipeline_args, image=image, generator=generator).images[0]
-            images.append(image)
+    for _ in range(args.num_validation_images):
+        with torch.autocast("cuda"):
+            image = pipeline(**pipeline_args, num_inference_steps=25, generator=generator).images[0]
+        images.append(image)
 
     for tracker in accelerator.trackers:
         if tracker.name == "tensorboard":
@@ -533,19 +525,6 @@ def parse_args(input_args=None):
     parser.add_argument(
         "--skip_save_text_encoder", action="store_true", required=False, help="Set to not save text encoder"
     )
-    parser.add_argument(
-        "--validation_images",
-        required=False,
-        default=None,
-        nargs="+",
-        help="Optional set of images to use for validation. Used when the target pipeline takes an initial image as input such as when training image variation or superresolution.",
-    )
-    parser.add_argument(
-        "--class_labels_conditioning",
-        required=False,
-        default=None,
-        help="The optional `class_label` conditioning to pass to the unet, available values are `timesteps`.",
-    )
 
     if input_args is not None:
         args = parser.parse_args(input_args)
@@ -772,12 +751,13 @@ def encode_prompt(text_encoder, input_ids, attention_mask, text_encoder_use_atte
 def main(args):
     logging_dir = Path(args.output_dir, args.logging_dir)
 
-    accelerator_project_config = ProjectConfiguration(project_dir=args.output_dir, logging_dir=logging_dir)
+    accelerator_project_config = ProjectConfiguration(total_limit=args.checkpoints_total_limit)
 
     accelerator = Accelerator(
         gradient_accumulation_steps=args.gradient_accumulation_steps,
         mixed_precision=args.mixed_precision,
         log_with=args.report_to,
+        logging_dir=logging_dir,
         project_config=accelerator_project_config,
     )
 
@@ -1091,8 +1071,8 @@ def main(args):
             unet, optimizer, train_dataloader, lr_scheduler
         )
 
-    # For mixed precision training we cast all non-trainable weigths (vae, non-lora text_encoder and non-lora unet) to half-precision
-    # as these weights are only used for inference, keeping weights in full precision is not required.
+    # For mixed precision training we cast the text_encoder and vae weights to half-precision
+    # as these models are only used for inference, keeping weights in full precision is not required.
     weight_dtype = torch.float32
     if accelerator.mixed_precision == "fp16":
         weight_dtype = torch.float16
@@ -1189,7 +1169,7 @@ def main(args):
                     )
                 else:
                     noise = torch.randn_like(model_input)
-                bsz, channels, height, width = model_input.shape
+                bsz = model_input.shape[0]
                 # Sample a random timestep for each image
                 timesteps = torch.randint(
                     0, noise_scheduler.config.num_train_timesteps, (bsz,), device=model_input.device
@@ -1211,18 +1191,8 @@ def main(args):
                         text_encoder_use_attention_mask=args.text_encoder_use_attention_mask,
                     )
 
-                if accelerator.unwrap_model(unet).config.in_channels == channels * 2:
-                    noisy_model_input = torch.cat([noisy_model_input, noisy_model_input], dim=1)
-
-                if args.class_labels_conditioning == "timesteps":
-                    class_labels = timesteps
-                else:
-                    class_labels = None
-
                 # Predict the noise residual
-                model_pred = unet(
-                    noisy_model_input, timesteps, encoder_hidden_states, class_labels=class_labels
-                ).sample
+                model_pred = unet(noisy_model_input, timesteps, encoder_hidden_states).sample
 
                 if model_pred.shape[1] == 6:
                     model_pred, _ = torch.chunk(model_pred, 2, dim=1)
@@ -1269,33 +1239,12 @@ def main(args):
                 global_step += 1
 
                 if accelerator.is_main_process:
+                    images = []
                     if global_step % args.checkpointing_steps == 0:
-                        # _before_ saving state, check if this save would set us over the `checkpoints_total_limit`
-                        if args.checkpoints_total_limit is not None:
-                            checkpoints = os.listdir(args.output_dir)
-                            checkpoints = [d for d in checkpoints if d.startswith("checkpoint")]
-                            checkpoints = sorted(checkpoints, key=lambda x: int(x.split("-")[1]))
-
-                            # before we save the new checkpoint, we need to have at _most_ `checkpoints_total_limit - 1` checkpoints
-                            if len(checkpoints) >= args.checkpoints_total_limit:
-                                num_to_remove = len(checkpoints) - args.checkpoints_total_limit + 1
-                                removing_checkpoints = checkpoints[0:num_to_remove]
-
-                                logger.info(
-                                    f"{len(checkpoints)} checkpoints already exist, removing {len(removing_checkpoints)} checkpoints"
-                                )
-                                logger.info(f"removing checkpoints: {', '.join(removing_checkpoints)}")
-
-                                for removing_checkpoint in removing_checkpoints:
-                                    removing_checkpoint = os.path.join(args.output_dir, removing_checkpoint)
-                                    shutil.rmtree(removing_checkpoint)
-
                         save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
                         accelerator.save_state(save_path)
                         logger.info(f"Saved state to {save_path}")
 
-                    images = []
-
                     if args.validation_prompt is not None and global_step % args.validation_steps == 0:
                         images = log_validation(
                             text_encoder,