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39 Commits

Author SHA1 Message Date
YiYi Xu 0b6dcf696a Merge branch 'main' into custom-block-local-code 2025-10-30 12:11:17 -10:00
Dhruv Nair df8dd77817 [Modular] Fix for custom block kwargs (#12561)
update
2025-10-31 00:14:24 +05:30
Pavle Padjin 9f3c0fdcd8 Avoiding graph break by changing the way we infer dtype in vae.decoder (#12512)
* Changing the way we infer dtype to avoid force evaluation of lazy tensors

* changing way to infer dtype to ensure type consistency

* more robust infering of dtype

* removing the upscale dtype entirely
2025-10-30 08:39:40 +05:30
galbria 84e16575e4 Bria fibo (#12545)
* Bria FIBO pipeline

* style fixs

* fix CR

* Refactor BriaFibo classes and update pipeline parameters

- Updated BriaFiboAttnProcessor and BriaFiboAttention classes to reflect changes from Flux equivalents.
- Modified the _unpack_latents method in BriaFiboPipeline to improve clarity.
- Increased the default max_sequence_length to 3000 and added a new optional parameter do_patching.
- Cleaned up test_pipeline_bria_fibo.py by removing unused imports and skipping unsupported tests.

* edit the docs of FIBO

* Remove unused BriaFibo imports and update CPU offload method in BriaFiboPipeline

* Refactor FIBO classes to BriaFibo naming convention

- Updated class names from FIBO to BriaFibo for consistency across the module.
- Modified instances of FIBOEmbedND, FIBOTimesteps, TextProjection, and TimestepProjEmbeddings to reflect the new naming.
- Ensured all references in the BriaFiboTransformer2DModel are updated accordingly.

* Add BriaFiboTransformer2DModel import to transformers module

* Remove unused BriaFibo imports from modular pipelines and add BriaFiboTransformer2DModel and BriaFiboPipeline classes to dummy objects for enhanced compatibility with torch and transformers.

* Update BriaFibo classes with copied documentation and fix import typo in pipeline module

- Added documentation comments indicating the source of copied code in BriaFiboTransformerBlock and _pack_latents methods.
- Corrected the import statement for BriaFiboPipeline in the pipelines module.

* Remove unused BriaFibo imports from __init__.py to streamline modular pipelines.

* Refactor documentation comments in BriaFibo classes to indicate inspiration from existing implementations

- Updated comments in BriaFiboAttnProcessor, BriaFiboAttention, and BriaFiboPipeline to reflect that the code is inspired by other modules rather than copied.
- Enhanced clarity on the origins of the methods to maintain proper attribution.

* change Inspired by to Based on

* add reference link and fix trailing whitespace

* Add BriaFiboTransformer2DModel documentation and update comments in BriaFibo classes

- Introduced a new documentation file for BriaFiboTransformer2DModel.
- Updated comments in BriaFiboAttnProcessor, BriaFiboAttention, and BriaFiboPipeline to clarify the origins of the code, indicating copied sources for better attribution.

---------

Co-authored-by: sayakpaul <spsayakpaul@gmail.com>
2025-10-28 16:27:48 +05:30
Sayak Paul 55d49d4379 [ci] don't run sana layerwise casting tests in CI. (#12551)
* don't run sana layerwise casting tests in CI.

* up
2025-10-28 13:29:51 +05:30
Meatfucker 40528e9ae7 Fix typos in kandinsky5 docs (#12552)
Update kandinsky5.md

Fix typos
2025-10-28 02:54:24 -03:00
Wang, Yi dc622a95d0 fix crash if tiling mode is enabled (#12521)
* fix crash in tiling mode is enabled

Signed-off-by: Wang, Yi A <yi.a.wang@intel.com>

* fmt

Signed-off-by: Wang, Yi A <yi.a.wang@intel.com>

---------

Signed-off-by: Wang, Yi A <yi.a.wang@intel.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-27 17:59:20 -10:00
Dhruv Nair ecfbc8f952 [Pipelines] Enable Wan VACE to run since single transformer (#12428)
* update

* update

* update

* update

* update
2025-10-28 09:21:31 +05:30
Sayak Paul df0e2a4f2c support latest few-step wan LoRA. (#12541)
* support latest few-step wan LoRA.

* up

* up
2025-10-28 08:55:24 +05:30
G.O.D 303efd2b8d Improve pos embed for Flux.1 inference on Ascend NPU (#12534)
improve pos embed for ascend npu

Co-authored-by: felix01.yu <felix01.yu@vipshop.com>
2025-10-27 16:55:36 -10:00
Lev Novitskiy 5afbcce176 Kandinsky 5 10 sec (NABLA suport) (#12520)
* add transformer pipeline first version

* updates

* fix 5sec generation

* rewrite Kandinsky5T2VPipeline to diffusers style

* add multiprompt support

* remove prints in pipeline

* add nabla attention

* Wrap Transformer in Diffusers style

* fix license

* fix prompt type

* add gradient checkpointing and peft support

* add usage example

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>

* remove unused imports

* add 10 second models support

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* remove no_grad and simplified prompt paddings

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* moved template to __init__

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* moved sdps inside processor

* remove oneline function

* remove reset_dtype methods

* Transformer: move all methods to forward

* separated prompt encoding

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* refactoring

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* refactoring acording to https://github.com/huggingface/diffusers/commit/acabbc0033d4b4933fc651766a4aa026db2e6dc1

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* Update src/diffusers/pipelines/kandinsky5/pipeline_kandinsky.py

Co-authored-by: YiYi Xu <yixu310@gmail.com>

* fixed

* style +copies

* Update src/diffusers/models/transformers/transformer_kandinsky.py

Co-authored-by: Charles <charles@huggingface.co>

* more

* Apply suggestions from code review

* add lora loader doc

* add compiled Nabla Attention

* all needed changes for 10 sec models are added!

* add docs

* Apply style fixes

* update docs

* add kandinsky5 to toctree

* add tests

* fix tests

* Apply style fixes

* update tests

---------

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Charles <charles@huggingface.co>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: github-actions[bot] <github-actions[bot]@users.noreply.github.com>
2025-10-28 07:47:18 +05:30
alirezafarashah 6d1a648602 Fix small inconsistency in output dimension of "_get_t5_prompt_embeds" function in sd3 pipeline (#12531)
* Fix small inconsistency in output dimension of t5 embeds when text_encoder_3 is None

* first commit

---------

Co-authored-by: Alireza Farashah <alireza.farashah@cn-g017.server.mila.quebec>
Co-authored-by: Alireza Farashah <alireza.farashah@login-2.server.mila.quebec>
2025-10-27 07:16:43 -10:00
Mikko Lauri 250f5cb53d Add AITER attention backend (#12549)
* add aiter attention backend

* Apply style fixes

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: github-actions[bot] <github-actions[bot]@users.noreply.github.com>
2025-10-27 20:25:02 +05:30
josephrocca dc6bd1511a Fix Chroma attention padding order and update docs to use lodestones/Chroma1-HD (#12508)
* [Fix] Move attention mask padding after T5 embedding

* [Fix] Move attention mask padding after T5 embedding

* Clean up whitespace in pipeline_chroma.py

Removed unnecessary blank lines for cleaner code.

* Fix

* Fix

* Update model to final Chroma1-HD checkpoint

* Update to Chroma1-HD

* Update model to Chroma1-HD

* Update model to Chroma1-HD

* Update Chroma model links to Chroma1-HD

* Add comment about padding/masking

* Fix checkpoint/repo references

* Apply style fixes

---------

Co-authored-by: github-actions[bot] <github-actions[bot]@users.noreply.github.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
2025-10-27 16:25:20 +05:30
Sayak Paul 500b9cf184 [chore] Move guiders experimental warning (#12543)
* move guiders experimental warning to init.

* up
2025-10-26 07:41:23 -10:00
Dhruv Nair d34b18c783 Deprecate Stable Cascade (#12537)
* update

* update

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-24 22:06:31 +05:30
kaixuanliu 7536f647e4 Loose the criteria tolerance appropriately for Intel XPU devices (#12460)
* Loose the criteria tolerance appropriately for Intel XPU devices

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>

* change back the atol value

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>

* use expectations

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>

* Update tests/pipelines/kandinsky2_2/test_kandinsky_controlnet.py

---------

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>
Co-authored-by: Ilyas Moutawwakil <57442720+IlyasMoutawwakil@users.noreply.github.com>
2025-10-24 12:18:15 +02:00
YiYi Xu a138d71ec1 HunyuanImage21 (#12333)
* add hunyuanimage2.1


---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-23 22:31:12 -10:00
Sayak Paul bc4039886d fix constants.py to user upper() (#12479) 2025-10-24 12:00:02 +05:30
Dhruv Nair 9c3b58dcf1 Handle deprecated transformer classes (#12517)
* update

* update

* update
2025-10-23 16:22:07 +05:30
Aishwarya Badlani 74b5fed434 Fix MPS compatibility in get_1d_sincos_pos_embed_from_grid #12432 (#12449)
* Fix MPS compatibility in get_1d_sincos_pos_embed_from_grid #12432

* Fix trailing whitespace in docstring

* Apply style fixes

---------

Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: github-actions[bot] <github-actions[bot]@users.noreply.github.com>
2025-10-23 16:18:07 +05:30
kaixuanliu 85eb505672 fix CI bug for kandinsky3_img2img case (#12474)
* fix CI bug for kandinsky3_img2img case

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>

* update code

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>

---------

Signed-off-by: Liu, Kaixuan <kaixuan.liu@intel.com>
2025-10-23 16:17:22 +05:30
Sayak Paul ccdd96ca52 [tests] Test attention backends (#12388)
* add a lightweight test suite for attention backends.

* up

* up

* Apply suggestions from code review

* formatting
2025-10-23 15:09:41 +05:30
Sayak Paul 4c723d8ec3 [CI] xfail the test_wuerstchen_prior test (#12530)
xfail the test_wuerstchen_prior test
2025-10-22 08:45:47 -10:00
YiYi Xu bec2d8eaea Fix: Add _skip_keys for AutoencoderKLWan (#12523)
add
2025-10-22 07:53:13 -10:00
Álvaro Somoza a0a51eb098 Kandinsky5 No cfg fix (#12527)
fix
2025-10-22 22:02:47 +05:30
Sayak Paul a5a0ccf86a [core] AutoencoderMixin to abstract common methods (#12473)
* up

* correct wording.

* up

* up

* up
2025-10-22 08:52:06 +05:30
David Bertoin dd07b19e27 Prx (#12525)
* rename photon to prx

* rename photon into prx

* Revert .gitignore to state before commit b7fb0fe9d6

* rename photon to prx

* rename photon into prx

* Revert .gitignore to state before commit b7fb0fe9d6

* make fix-copies
2025-10-21 17:09:22 -07:00
vb 57636ad4f4 purge HF_HUB_ENABLE_HF_TRANSFER; promote Xet (#12497)
* purge HF_HUB_ENABLE_HF_TRANSFER; promote Xet

* purge HF_HUB_ENABLE_HF_TRANSFER; promote Xet x2

* restrict docker build test to the ones we actually use in CI.

---------

Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-22 00:59:20 +05:30
David Bertoin cefc2cf82d Add Photon model and pipeline support (#12456)
* Add Photon model and pipeline support

This commit adds support for the Photon image generation model:
- PhotonTransformer2DModel: Core transformer architecture
- PhotonPipeline: Text-to-image generation pipeline
- Attention processor updates for Photon-specific attention mechanism
- Conversion script for loading Photon checkpoints
- Documentation and tests

* just store the T5Gemma encoder

* enhance_vae_properties if vae is provided only

* remove autocast for text encoder forwad

* BF16 example

* conditioned CFG

* remove enhance vae and use vae.config directly when possible

* move PhotonAttnProcessor2_0 in transformer_photon

* remove einops dependency and now inherits from AttentionMixin

* unify the structure of the forward block

* update doc

* update doc

* fix T5Gemma loading from hub

* fix timestep shift

* remove lora support from doc

* Rename EmbedND for PhotoEmbedND

* remove modulation dataclass

* put _attn_forward and _ffn_forward logic in PhotonBlock's forward

* renam LastLayer for FinalLayer

* remove lora related code

* rename vae_spatial_compression_ratio for vae_scale_factor

* support prompt_embeds in call

* move xattention conditionning out computation out of the denoising loop

* add negative prompts

* Use _import_structure for lazy loading

* make quality + style

* add pipeline test + corresponding fixes

* utility function that determines the default resolution given the VAE

* Refactor PhotonAttention to match Flux pattern

* built-in RMSNorm

* Revert accidental .gitignore change

* parameter names match the standard diffusers conventions

* renaming and remove unecessary attributes setting

* Update docs/source/en/api/pipelines/photon.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* quantization example

* added doc to toctree

* Update docs/source/en/api/pipelines/photon.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/photon.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* Update docs/source/en/api/pipelines/photon.md

Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>

* use dispatch_attention_fn for multiple attention backend support

* naming changes

* make fix copy

* Update docs/source/en/api/pipelines/photon.md

Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>

* Add PhotonTransformer2DModel to TYPE_CHECKING imports

* make fix-copies

* Use Tuple instead of tuple

Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>

* restrict the version of transformers

Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>

* Update tests/pipelines/photon/test_pipeline_photon.py

Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>

* Update tests/pipelines/photon/test_pipeline_photon.py

Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>

* change | for Optional

* fix nits.

* use typing Dict

---------

Co-authored-by: davidb <davidb@worker-10.soperator-worker-svc.soperator.svc.cluster.local>
Co-authored-by: David Briand <david@photoroom.com>
Co-authored-by: davidb <davidb@worker-8.soperator-worker-svc.soperator.svc.cluster.local>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: dg845 <58458699+dg845@users.noreply.github.com>
Co-authored-by: sayakpaul <spsayakpaul@gmail.com>
2025-10-21 20:55:55 +05:30
Sayak Paul b3e56e71fb styling issues. (#12522) 2025-10-21 20:04:54 +05:30
Steven Liu 5b5fa49a89 [docs] Organize toctree by modality (#12514)
* reorganize

* fix

---------

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>
2025-10-21 10:18:54 +05:30
Fei Xie decfa3c9e1 Fix: Use incorrect temporary variable key when replacing adapter name… (#12502)
Fix: Use incorrect temporary variable key when replacing adapter name in state dict within load_lora_adapter function

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-20 15:45:37 -10:00
Dhruv Nair 48305755bf Raise warning instead of error when imports are missing for custom code (#12513)
update
2025-10-20 07:02:23 -10:00
dg845 7853bfbed7 Remove Qwen Image Redundant RoPE Cache (#12452)
Refactor QwenEmbedRope to only use the LRU cache for RoPE caching
2025-10-19 18:41:58 -07:00
Lev Novitskiy 23ebbb4bc8 Kandinsky 5 is finally in Diffusers! (#12478)
* add kandinsky5 transformer pipeline first version

---------

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Charles <charles@huggingface.co>
2025-10-17 18:34:30 -10:00
Ali Imran 1b456bd5d5 docs: cleanup of runway model (#12503)
* cleanup of runway model

* quality fixes
2025-10-17 14:10:50 -07:00
Sayak Paul af769881d3 [tests] introduce VAETesterMixin to consolidate tests for slicing and tiling (#12374)
* up

* up

* up

* up

* up

* u[

* up

* up

* up
2025-10-17 12:02:29 +05:30
DN6 351b2f172a update 2025-09-24 14:33:38 +05:30
241 changed files with 14753 additions and 1245 deletions
+1 -1
View File
@@ -7,7 +7,7 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
+29 -8
View File
@@ -42,18 +42,39 @@ jobs:
CHANGED_FILES: ${{ steps.file_changes.outputs.all }}
run: |
echo "$CHANGED_FILES"
for FILE in $CHANGED_FILES; do
ALLOWED_IMAGES=(
diffusers-pytorch-cpu
diffusers-pytorch-cuda
diffusers-pytorch-xformers-cuda
diffusers-pytorch-minimum-cuda
diffusers-doc-builder
)
declare -A IMAGES_TO_BUILD=()
for FILE in $CHANGED_FILES; do
# skip anything that isn't still on disk
if [[ ! -f "$FILE" ]]; then
if [[ ! -e "$FILE" ]]; then
echo "Skipping removed file $FILE"
continue
fi
if [[ "$FILE" == docker/*Dockerfile ]]; then
DOCKER_PATH="${FILE%/Dockerfile}"
DOCKER_TAG=$(basename "$DOCKER_PATH")
echo "Building Docker image for $DOCKER_TAG"
docker build -t "$DOCKER_TAG" "$DOCKER_PATH"
fi
for IMAGE in "${ALLOWED_IMAGES[@]}"; do
if [[ "$FILE" == docker/${IMAGE}/* ]]; then
IMAGES_TO_BUILD["$IMAGE"]=1
fi
done
done
if [[ ${#IMAGES_TO_BUILD[@]} -eq 0 ]]; then
echo "No relevant Docker changes detected."
exit 0
fi
for IMAGE in "${!IMAGES_TO_BUILD[@]}"; do
DOCKER_PATH="docker/${IMAGE}"
echo "Building Docker image for $IMAGE"
docker build -t "$IMAGE" "$DOCKER_PATH"
done
if: steps.file_changes.outputs.all != ''
+1 -1
View File
@@ -7,7 +7,7 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
+1 -1
View File
@@ -26,7 +26,7 @@ concurrency:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
+1 -1
View File
@@ -22,7 +22,7 @@ concurrency:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
+1 -1
View File
@@ -24,7 +24,7 @@ env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 1000000000 # set high cutoff so that only always-test pipelines run
+1 -1
View File
@@ -14,7 +14,7 @@ env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 50000
+1 -1
View File
@@ -18,7 +18,7 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
PYTEST_TIMEOUT: 600
RUN_SLOW: no
+1 -1
View File
@@ -8,7 +8,7 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_XET_HIGH_PERFORMANCE: 1
PYTEST_TIMEOUT: 600
RUN_SLOW: no
+1 -1
View File
@@ -171,7 +171,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
<tr style="border-top: 2px solid black">
<td>Text-guided Image Inpainting</td>
<td><a href="https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/inpaint">Stable Diffusion Inpainting</a></td>
<td><a href="https://huggingface.co/runwayml/stable-diffusion-inpainting"> runwayml/stable-diffusion-inpainting </a></td>
<td><a href="https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting"> stable-diffusion-v1-5/stable-diffusion-inpainting </a></td>
</tr>
<tr style="border-top: 2px solid black">
<td>Image Variation</td>
+1 -1
View File
@@ -33,7 +33,7 @@ RUN uv pip install --no-cache-dir "git+https://github.com/huggingface/diffusers.
RUN uv pip install --no-cache-dir \
accelerate \
numpy==1.26.4 \
hf_transfer \
hf_xet \
setuptools==69.5.1 \
bitsandbytes \
torchao \
+1 -1
View File
@@ -44,6 +44,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
scipy \
tensorboard \
transformers \
hf_transfer
hf_xet
CMD ["/bin/bash"]
+2 -3
View File
@@ -38,13 +38,12 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
hf_xet \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]
+1 -1
View File
@@ -31,7 +31,7 @@ RUN uv pip install --no-cache-dir "git+https://github.com/huggingface/diffusers.
RUN uv pip install --no-cache-dir \
accelerate \
numpy==1.26.4 \
hf_transfer
hf_xet
RUN apt-get clean && rm -rf /var/lib/apt/lists/* && apt-get autoremove && apt-get autoclean
+1 -1
View File
@@ -44,6 +44,6 @@ RUN uv pip install --no-cache-dir \
accelerate \
numpy==1.26.4 \
pytorch-lightning \
hf_transfer
hf_xet
CMD ["/bin/bash"]
@@ -47,6 +47,6 @@ RUN uv pip install --no-cache-dir \
accelerate \
numpy==1.26.4 \
pytorch-lightning \
hf_transfer
hf_xet
CMD ["/bin/bash"]
@@ -44,7 +44,7 @@ RUN uv pip install --no-cache-dir \
accelerate \
numpy==1.26.4 \
pytorch-lightning \
hf_transfer \
hf_xet \
xformers
CMD ["/bin/bash"]
+274 -263
View File
@@ -1,5 +1,4 @@
- title: Get started
sections:
- sections:
- local: index
title: Diffusers
- local: installation
@@ -8,9 +7,8 @@
title: Quickstart
- local: stable_diffusion
title: Basic performance
- title: Pipelines
isExpanded: false
title: Get started
- isExpanded: false
sections:
- local: using-diffusers/loading
title: DiffusionPipeline
@@ -28,9 +26,8 @@
title: Model formats
- local: using-diffusers/push_to_hub
title: Sharing pipelines and models
- title: Adapters
isExpanded: false
title: Pipelines
- isExpanded: false
sections:
- local: tutorials/using_peft_for_inference
title: LoRA
@@ -44,9 +41,8 @@
title: DreamBooth
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- title: Inference
isExpanded: false
title: Adapters
- isExpanded: false
sections:
- local: using-diffusers/weighted_prompts
title: Prompting
@@ -56,9 +52,8 @@
title: Batch inference
- local: training/distributed_inference
title: Distributed inference
- title: Inference optimization
isExpanded: false
title: Inference
- isExpanded: false
sections:
- local: optimization/fp16
title: Accelerate inference
@@ -70,8 +65,7 @@
title: Reduce memory usage
- local: optimization/speed-memory-optims
title: Compiling and offloading quantized models
- title: Community optimizations
sections:
- sections:
- local: optimization/pruna
title: Pruna
- local: optimization/xformers
@@ -90,9 +84,9 @@
title: ParaAttention
- local: using-diffusers/image_quality
title: FreeU
- title: Hybrid Inference
isExpanded: false
title: Community optimizations
title: Inference optimization
- isExpanded: false
sections:
- local: hybrid_inference/overview
title: Overview
@@ -102,9 +96,8 @@
title: VAE Encode
- local: hybrid_inference/api_reference
title: API Reference
- title: Modular Diffusers
isExpanded: false
title: Hybrid Inference
- isExpanded: false
sections:
- local: modular_diffusers/overview
title: Overview
@@ -126,9 +119,8 @@
title: ComponentsManager
- local: modular_diffusers/guiders
title: Guiders
- title: Training
isExpanded: false
title: Modular Diffusers
- isExpanded: false
sections:
- local: training/overview
title: Overview
@@ -138,8 +130,7 @@
title: Adapt a model to a new task
- local: tutorials/basic_training
title: Train a diffusion model
- title: Models
sections:
- sections:
- local: training/unconditional_training
title: Unconditional image generation
- local: training/text2image
@@ -158,8 +149,8 @@
title: InstructPix2Pix
- local: training/cogvideox
title: CogVideoX
- title: Methods
sections:
title: Models
- sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
@@ -172,9 +163,9 @@
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
- title: Quantization
isExpanded: false
title: Methods
title: Training
- isExpanded: false
sections:
- local: quantization/overview
title: Getting started
@@ -188,9 +179,8 @@
title: quanto
- local: quantization/modelopt
title: NVIDIA ModelOpt
- title: Model accelerators and hardware
isExpanded: false
title: Quantization
- isExpanded: false
sections:
- local: optimization/onnx
title: ONNX
@@ -204,9 +194,8 @@
title: Intel Gaudi
- local: optimization/neuron
title: AWS Neuron
- title: Specific pipeline examples
isExpanded: false
title: Model accelerators and hardware
- isExpanded: false
sections:
- local: using-diffusers/consisid
title: ConsisID
@@ -232,12 +221,10 @@
title: Stable Video Diffusion
- local: using-diffusers/marigold_usage
title: Marigold Computer Vision
- title: Resources
isExpanded: false
title: Specific pipeline examples
- isExpanded: false
sections:
- title: Task recipes
sections:
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: using-diffusers/conditional_image_generation
@@ -252,6 +239,7 @@
title: Video generation
- local: using-diffusers/depth2img
title: Depth-to-image
title: Task recipes
- local: using-diffusers/write_own_pipeline
title: Understanding pipelines, models and schedulers
- local: community_projects
@@ -266,12 +254,10 @@
title: Diffusers' Ethical Guidelines
- local: conceptual/evaluation
title: Evaluating Diffusion Models
- title: API
isExpanded: false
title: Resources
- isExpanded: false
sections:
- title: Main Classes
sections:
- sections:
- local: api/configuration
title: Configuration
- local: api/logging
@@ -282,8 +268,8 @@
title: Quantization
- local: api/parallel
title: Parallel inference
- title: Modular
sections:
title: Main Classes
- sections:
- local: api/modular_diffusers/pipeline
title: Pipeline
- local: api/modular_diffusers/pipeline_blocks
@@ -294,8 +280,8 @@
title: Components and configs
- local: api/modular_diffusers/guiders
title: Guiders
- title: Loaders
sections:
title: Modular
- sections:
- local: api/loaders/ip_adapter
title: IP-Adapter
- local: api/loaders/lora
@@ -310,14 +296,13 @@
title: SD3Transformer2D
- local: api/loaders/peft
title: PEFT
- title: Models
sections:
title: Loaders
- sections:
- local: api/models/overview
title: Overview
- local: api/models/auto_model
title: AutoModel
- title: ControlNets
sections:
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_union
@@ -332,12 +317,14 @@
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- title: Transformers
sections:
title: ControlNets
- sections:
- local: api/models/allegro_transformer3d
title: AllegroTransformer3DModel
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/transformer_bria_fibo
title: BriaFiboTransformer2DModel
- local: api/models/bria_transformer
title: BriaTransformer2DModel
- local: api/models/chroma_transformer
@@ -362,6 +349,8 @@
title: HiDreamImageTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/hunyuanimage_transformer_2d
title: HunyuanImageTransformer2DModel
- local: api/models/hunyuan_video_transformer_3d
title: HunyuanVideoTransformer3DModel
- local: api/models/latte_transformer3d
@@ -396,8 +385,8 @@
title: TransformerTemporalModel
- local: api/models/wan_transformer_3d
title: WanTransformer3DModel
- title: UNets
sections:
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/unet
@@ -412,8 +401,8 @@
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
- title: VAEs
sections:
title: UNets
- sections:
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
@@ -426,6 +415,10 @@
title: AutoencoderKLCogVideoX
- local: api/models/autoencoderkl_cosmos
title: AutoencoderKLCosmos
- local: api/models/autoencoder_kl_hunyuanimage
title: AutoencoderKLHunyuanImage
- local: api/models/autoencoder_kl_hunyuanimage_refiner
title: AutoencoderKLHunyuanImageRefiner
- local: api/models/autoencoder_kl_hunyuan_video
title: AutoencoderKLHunyuanVideo
- local: api/models/autoencoderkl_ltx_video
@@ -446,210 +439,226 @@
title: Tiny AutoEncoder
- local: api/models/vq
title: VQModel
- title: Pipelines
sections:
title: VAEs
title: Models
- sections:
- local: api/pipelines/overview
title: Overview
- local: api/pipelines/allegro
title: Allegro
- local: api/pipelines/amused
title: aMUSEd
- local: api/pipelines/animatediff
title: AnimateDiff
- local: api/pipelines/attend_and_excite
title: Attend-and-Excite
- local: api/pipelines/audioldm
title: AudioLDM
- local: api/pipelines/audioldm2
title: AudioLDM 2
- local: api/pipelines/aura_flow
title: AuraFlow
- sections:
- local: api/pipelines/audioldm
title: AudioLDM
- local: api/pipelines/audioldm2
title: AudioLDM 2
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/stable_audio
title: Stable Audio
title: Audio
- local: api/pipelines/auto_pipeline
title: AutoPipeline
- local: api/pipelines/blip_diffusion
title: BLIP-Diffusion
- local: api/pipelines/bria_3_2
title: Bria 3.2
- local: api/pipelines/chroma
title: Chroma
- local: api/pipelines/cogvideox
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_flux
title: ControlNet with Flux.1
- local: api/pipelines/controlnet_hunyuandit
title: ControlNet with Hunyuan-DiT
- local: api/pipelines/controlnet_sd3
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/cosmos
title: Cosmos
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
title: DDIM
- local: api/pipelines/ddpm
title: DDPM
- local: api/pipelines/deepfloyd_if
title: DeepFloyd IF
- local: api/pipelines/diffedit
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/framepack
title: Framepack
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/hunyuan_video
title: HunyuanVideo
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/pix2pix
title: InstructPix2Pix
- local: api/pipelines/kandinsky
title: Kandinsky 2.1
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/kandinsky3
title: Kandinsky 3
- local: api/pipelines/kolors
title: Kolors
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/latte
title: Latte
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/mochi
title: Mochi
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
title: Paint by Example
- local: api/pipelines/pia
title: Personalized Image Animator (PIA)
- local: api/pipelines/pixart
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/qwenimage
title: QwenImage
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/skyreels_v2
title: SkyReels-V2
- local: api/pipelines/stable_audio
title: Stable Audio
- local: api/pipelines/stable_cascade
title: Stable Cascade
- title: Stable Diffusion
sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-image
- sections:
- local: api/pipelines/amused
title: aMUSEd
- local: api/pipelines/animatediff
title: AnimateDiff
- local: api/pipelines/attend_and_excite
title: Attend-and-Excite
- local: api/pipelines/aura_flow
title: AuraFlow
- local: api/pipelines/blip_diffusion
title: BLIP-Diffusion
- local: api/pipelines/bria_3_2
title: Bria 3.2
- local: api/pipelines/bria_fibo
title: Bria Fibo
- local: api/pipelines/chroma
title: Chroma
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_flux
title: ControlNet with Flux.1
- local: api/pipelines/controlnet_hunyuandit
title: ControlNet with Hunyuan-DiT
- local: api/pipelines/controlnet_sd3
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/cosmos
title: Cosmos
- local: api/pipelines/ddim
title: DDIM
- local: api/pipelines/ddpm
title: DDPM
- local: api/pipelines/deepfloyd_if
title: DeepFloyd IF
- local: api/pipelines/diffedit
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/pix2pix
title: InstructPix2Pix
- local: api/pipelines/kandinsky
title: Kandinsky 2.1
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/kandinsky3
title: Kandinsky 3
- local: api/pipelines/kandinsky5
title: Kandinsky 5
- local: api/pipelines/kolors
title: Kolors
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
title: Paint by Example
- local: api/pipelines/pixart
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/prx
title: PRX
- local: api/pipelines/qwenimage
title: QwenImage
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/stable_cascade
title: Stable Cascade
- sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-image
- local: api/pipelines/stable_diffusion/inpaint
title: Inpainting
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D
Upscaler
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
title: Stable Diffusion
- local: api/pipelines/stable_unclip
title: Stable unCLIP
- local: api/pipelines/unclip
title: unCLIP
- local: api/pipelines/unidiffuser
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/visualcloze
title: VisualCloze
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Image
- sections:
- local: api/pipelines/allegro
title: Allegro
- local: api/pipelines/cogvideox
title: CogVideoX
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/framepack
title: Framepack
- local: api/pipelines/hunyuanimage21
title: HunyuanImage2.1
- local: api/pipelines/hunyuan_video
title: HunyuanVideo
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/latte
title: Latte
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/mochi
title: Mochi
- local: api/pipelines/pia
title: Personalized Image Animator (PIA)
- local: api/pipelines/skyreels_v2
title: SkyReels-V2
- local: api/pipelines/stable_diffusion/svd
title: Image-to-video
- local: api/pipelines/stable_diffusion/inpaint
title: Inpainting
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
- local: api/pipelines/stable_unclip
title: Stable unCLIP
- local: api/pipelines/text_to_video
title: Text-to-video
- local: api/pipelines/text_to_video_zero
title: Text2Video-Zero
- local: api/pipelines/unclip
title: unCLIP
- local: api/pipelines/unidiffuser
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/visualcloze
title: VisualCloze
- local: api/pipelines/wan
title: Wan
- local: api/pipelines/wuerstchen
title: Wuerstchen
- title: Schedulers
sections:
title: Stable Video Diffusion
- local: api/pipelines/text_to_video
title: Text-to-video
- local: api/pipelines/text_to_video_zero
title: Text2Video-Zero
- local: api/pipelines/wan
title: Wan
title: Video
title: Pipelines
- sections:
- local: api/schedulers/overview
title: Overview
- local: api/schedulers/cm_stochastic_iterative
@@ -718,8 +727,8 @@
title: UniPCMultistepScheduler
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
- title: Internal classes
sections:
title: Schedulers
- sections:
- local: api/internal_classes_overview
title: Overview
- local: api/attnprocessor
@@ -736,3 +745,5 @@
title: VAE Image Processor
- local: api/video_processor
title: Video Processor
title: Internal classes
title: API
+3
View File
@@ -107,6 +107,9 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
[[autodoc]] loaders.lora_pipeline.QwenImageLoraLoaderMixin
## KandinskyLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.KandinskyLoraLoaderMixin
## LoraBaseMixin
[[autodoc]] loaders.lora_base.LoraBaseMixin
@@ -39,7 +39,7 @@ mask_url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images
original_image = load_image(img_url).resize((512, 512))
mask_image = load_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained("runwayml/stable-diffusion-inpainting")
pipe = StableDiffusionInpaintPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting")
pipe.vae = AsymmetricAutoencoderKL.from_pretrained("cross-attention/asymmetric-autoencoder-kl-x-1-5")
pipe.to("cuda")
@@ -0,0 +1,32 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLHunyuanImage
The 2D variational autoencoder (VAE) model with KL loss used in [HunyuanImage2.1].
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLHunyuanImage
vae = AutoencoderKLHunyuanImage.from_pretrained("hunyuanvideo-community/HunyuanImage-2.1-Diffusers", subfolder="vae", torch_dtype=torch.bfloat16)
```
## AutoencoderKLHunyuanImage
[[autodoc]] AutoencoderKLHunyuanImage
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput
@@ -0,0 +1,32 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLHunyuanImageRefiner
The 3D variational autoencoder (VAE) model with KL loss used in [HunyuanImage2.1](https://github.com/Tencent-Hunyuan/HunyuanImage-2.1) for its refiner pipeline.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLHunyuanImageRefiner
vae = AutoencoderKLHunyuanImageRefiner.from_pretrained("hunyuanvideo-community/HunyuanImage-2.1-Refiner-Diffusers", subfolder="vae", torch_dtype=torch.bfloat16)
```
## AutoencoderKLHunyuanImageRefiner
[[autodoc]] AutoencoderKLHunyuanImageRefiner
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput
@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# ChromaTransformer2DModel
A modified flux Transformer model from [Chroma](https://huggingface.co/lodestones/Chroma)
A modified flux Transformer model from [Chroma](https://huggingface.co/lodestones/Chroma1-HD)
## ChromaTransformer2DModel
@@ -0,0 +1,30 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# HunyuanImageTransformer2DModel
A Diffusion Transformer model for [HunyuanImage2.1](https://github.com/Tencent-Hunyuan/HunyuanImage-2.1).
The model can be loaded with the following code snippet.
```python
from diffusers import HunyuanImageTransformer2DModel
transformer = HunyuanImageTransformer2DModel.from_pretrained("hunyuanvideo-community/HunyuanImage-2.1-Diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## HunyuanImageTransformer2DModel
[[autodoc]] HunyuanImageTransformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
@@ -0,0 +1,19 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# BriaFiboTransformer2DModel
A modified flux Transformer model from [Bria](https://huggingface.co/briaai/FIBO)
## BriaFiboTransformer2DModel
[[autodoc]] BriaFiboTransformer2DModel
+45
View File
@@ -0,0 +1,45 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Bria Fibo
Text-to-image models have mastered imagination - but not control. FIBO changes that.
FIBO is trained on structured JSON captions up to 1,000+ words and designed to understand and control different visual parameters such as lighting, composition, color, and camera settings, enabling precise and reproducible outputs.
With only 8 billion parameters, FIBO provides a new level of image quality, prompt adherence and proffesional control.
FIBO is trained exclusively on a structured prompt and will not work with freeform text prompts.
you can use the [FIBO-VLM-prompt-to-JSON](https://huggingface.co/briaai/FIBO-VLM-prompt-to-JSON) model or the [FIBO-gemini-prompt-to-JSON](https://huggingface.co/briaai/FIBO-gemini-prompt-to-JSON) to convert your freeform text prompt to a structured JSON prompt.
its not recommended to use freeform text prompts directly with FIBO, as it will not produce the best results.
you can learn more about FIBO in [Bria Fibo Hugging Face page](https://huggingface.co/briaai/FIBO).
## Usage
_As the model is gated, before using it with diffusers you first need to go to the [Bria Fibo Hugging Face page](https://huggingface.co/briaai/FIBO), fill in the form and accept the gate. Once you are in, you need to login so that your system knows youve accepted the gate._
Use the command below to log in:
```bash
hf auth login
```
## BriaPipeline
[[autodoc]] BriaPipeline
- all
- __call__
+7 -6
View File
@@ -19,20 +19,21 @@ specific language governing permissions and limitations under the License.
Chroma is a text to image generation model based on Flux.
Original model checkpoints for Chroma can be found [here](https://huggingface.co/lodestones/Chroma).
Original model checkpoints for Chroma can be found here:
* High-resolution finetune: [lodestones/Chroma1-HD](https://huggingface.co/lodestones/Chroma1-HD)
* Base model: [lodestones/Chroma1-Base](https://huggingface.co/lodestones/Chroma1-Base)
* Original repo with progress checkpoints: [lodestones/Chroma](https://huggingface.co/lodestones/Chroma) (loading this repo with `from_pretrained` will load a Diffusers-compatible version of the `unlocked-v37` checkpoint)
> [!TIP]
> Chroma can use all the same optimizations as Flux.
## Inference
The Diffusers version of Chroma is based on the [`unlocked-v37`](https://huggingface.co/lodestones/Chroma/blob/main/chroma-unlocked-v37.safetensors) version of the original model, which is available in the [Chroma repository](https://huggingface.co/lodestones/Chroma).
```python
import torch
from diffusers import ChromaPipeline
pipe = ChromaPipeline.from_pretrained("lodestones/Chroma", torch_dtype=torch.bfloat16)
pipe = ChromaPipeline.from_pretrained("lodestones/Chroma1-HD", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = [
@@ -63,10 +64,10 @@ Then run the following example
import torch
from diffusers import ChromaTransformer2DModel, ChromaPipeline
model_id = "lodestones/Chroma"
model_id = "lodestones/Chroma1-HD"
dtype = torch.bfloat16
transformer = ChromaTransformer2DModel.from_single_file("https://huggingface.co/lodestones/Chroma/blob/main/chroma-unlocked-v37.safetensors", torch_dtype=dtype)
transformer = ChromaTransformer2DModel.from_single_file("https://huggingface.co/lodestones/Chroma1-HD/blob/main/Chroma1-HD.safetensors", torch_dtype=dtype)
pipe = ChromaPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=dtype)
pipe.enable_model_cpu_offload()
@@ -0,0 +1,152 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# HunyuanImage2.1
HunyuanImage-2.1 is a 17B text-to-image model that is capable of generating 2K (2048 x 2048) resolution images
HunyuanImage-2.1 comes in the following variants:
| model type | model id |
|:----------:|:--------:|
| HunyuanImage-2.1 | [hunyuanvideo-community/HunyuanImage-2.1-Diffusers](https://huggingface.co/hunyuanvideo-community/HunyuanImage-2.1-Diffusers) |
| HunyuanImage-2.1-Distilled | [hunyuanvideo-community/HunyuanImage-2.1-Distilled-Diffusers](https://huggingface.co/hunyuanvideo-community/HunyuanImage-2.1-Distilled-Diffusers) |
| HunyuanImage-2.1-Refiner | [hunyuanvideo-community/HunyuanImage-2.1-Refiner-Diffusers](https://huggingface.co/hunyuanvideo-community/HunyuanImage-2.1-Refiner-Diffusers) |
> [!TIP]
> [Caching](../../optimization/cache) may also speed up inference by storing and reusing intermediate outputs.
## HunyuanImage-2.1
HunyuanImage-2.1 applies [Adaptive Projected Guidance (APG)](https://huggingface.co/papers/2410.02416) combined with Classifier-Free Guidance (CFG) in the denoising loop. `HunyuanImagePipeline` has a `guider` component (read more about [Guider](../modular_diffusers/guiders.md)) and does not take a `guidance_scale` parameter at runtime. To change guider-related parameters, e.g., `guidance_scale`, you can update the `guider` configuration instead.
```python
import torch
from diffusers import HunyuanImagePipeline
pipe = HunyuanImagePipeline.from_pretrained(
"hunyuanvideo-community/HunyuanImage-2.1-Diffusers",
torch_dtype=torch.bfloat16
)
pipe = pipe.to("cuda")
```
You can inspect the `guider` object:
```py
>>> pipe.guider
AdaptiveProjectedMixGuidance {
"_class_name": "AdaptiveProjectedMixGuidance",
"_diffusers_version": "0.36.0.dev0",
"adaptive_projected_guidance_momentum": -0.5,
"adaptive_projected_guidance_rescale": 10.0,
"adaptive_projected_guidance_scale": 10.0,
"adaptive_projected_guidance_start_step": 5,
"enabled": true,
"eta": 0.0,
"guidance_rescale": 0.0,
"guidance_scale": 3.5,
"start": 0.0,
"stop": 1.0,
"use_original_formulation": false
}
State:
step: None
num_inference_steps: None
timestep: None
count_prepared: 0
enabled: True
num_conditions: 2
momentum_buffer: None
is_apg_enabled: False
is_cfg_enabled: True
```
To update the guider with a different configuration, use the `new()` method. For example, to generate an image with `guidance_scale=5.0` while keeping all other default guidance parameters:
```py
import torch
from diffusers import HunyuanImagePipeline
pipe = HunyuanImagePipeline.from_pretrained(
"hunyuanvideo-community/HunyuanImage-2.1-Diffusers",
torch_dtype=torch.bfloat16
)
pipe = pipe.to("cuda")
# Update the guider configuration
pipe.guider = pipe.guider.new(guidance_scale=5.0)
prompt = (
"A cute, cartoon-style anthropomorphic penguin plush toy with fluffy fur, standing in a painting studio, "
"wearing a red knitted scarf and a red beret with the word 'Tencent' on it, holding a paintbrush with a "
"focused expression as it paints an oil painting of the Mona Lisa, rendered in a photorealistic photographic style."
)
image = pipe(
prompt=prompt,
num_inference_steps=50,
height=2048,
width=2048,
).images[0]
image.save("image.png")
```
## HunyuanImage-2.1-Distilled
use `distilled_guidance_scale` with the guidance-distilled checkpoint,
```py
import torch
from diffusers import HunyuanImagePipeline
pipe = HunyuanImagePipeline.from_pretrained("hunyuanvideo-community/HunyuanImage-2.1-Distilled-Diffusers", torch_dtype=torch.bfloat16)
pipe = pipe.to("cuda")
prompt = (
"A cute, cartoon-style anthropomorphic penguin plush toy with fluffy fur, standing in a painting studio, "
"wearing a red knitted scarf and a red beret with the word 'Tencent' on it, holding a paintbrush with a "
"focused expression as it paints an oil painting of the Mona Lisa, rendered in a photorealistic photographic style."
)
out = pipe(
prompt,
num_inference_steps=8,
distilled_guidance_scale=3.25,
height=2048,
width=2048,
generator=generator,
).images[0]
```
## HunyuanImagePipeline
[[autodoc]] HunyuanImagePipeline
- all
- __call__
## HunyuanImageRefinerPipeline
[[autodoc]] HunyuanImageRefinerPipeline
- all
- __call__
## HunyuanImagePipelineOutput
[[autodoc]] pipelines.hunyuan_image.pipeline_output.HunyuanImagePipelineOutput
+149
View File
@@ -0,0 +1,149 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Kandinsky 5.0
Kandinsky 5.0 is created by the Kandinsky team: Alexey Letunovskiy, Maria Kovaleva, Ivan Kirillov, Lev Novitskiy, Denis Koposov, Dmitrii Mikhailov, Anna Averchenkova, Andrey Shutkin, Julia Agafonova, Olga Kim, Anastasiia Kargapoltseva, Nikita Kiselev, Anna Dmitrienko, Anastasia Maltseva, Kirill Chernyshev, Ilia Vasiliev, Viacheslav Vasilev, Vladimir Polovnikov, Yury Kolabushin, Alexander Belykh, Mikhail Mamaev, Anastasia Aliaskina, Tatiana Nikulina, Polina Gavrilova, Vladimir Arkhipkin, Vladimir Korviakov, Nikolai Gerasimenko, Denis Parkhomenko, Denis Dimitrov
Kandinsky 5.0 is a family of diffusion models for Video & Image generation. Kandinsky 5.0 T2V Lite is a lightweight video generation model (2B parameters) that ranks #1 among open-source models in its class. It outperforms larger models and offers the best understanding of Russian concepts in the open-source ecosystem.
The model introduces several key innovations:
- **Latent diffusion pipeline** with **Flow Matching** for improved training stability
- **Diffusion Transformer (DiT)** as the main generative backbone with cross-attention to text embeddings
- Dual text encoding using **Qwen2.5-VL** and **CLIP** for comprehensive text understanding
- **HunyuanVideo 3D VAE** for efficient video encoding and decoding
- **Sparse attention mechanisms** (NABLA) for efficient long-sequence processing
The original codebase can be found at [ai-forever/Kandinsky-5](https://github.com/ai-forever/Kandinsky-5).
> [!TIP]
> Check out the [AI Forever](https://huggingface.co/ai-forever) organization on the Hub for the official model checkpoints for text-to-video generation, including pretrained, SFT, no-CFG, and distilled variants.
## Available Models
Kandinsky 5.0 T2V Lite comes in several variants optimized for different use cases:
| model_id | Description | Use Cases |
|------------|-------------|-----------|
| **ai-forever/Kandinsky-5.0-T2V-Lite-sft-5s-Diffusers** | 5 second Supervised Fine-Tuned model | Highest generation quality |
| **ai-forever/Kandinsky-5.0-T2V-Lite-sft-10s-Diffusers** | 10 second Supervised Fine-Tuned model | Highest generation quality |
| **ai-forever/Kandinsky-5.0-T2V-Lite-nocfg-5s-Diffusers** | 5 second Classifier-Free Guidance distilled | 2× faster inference |
| **ai-forever/Kandinsky-5.0-T2V-Lite-nocfg-10s-Diffusers** | 10 second Classifier-Free Guidance distilled | 2× faster inference |
| **ai-forever/Kandinsky-5.0-T2V-Lite-distilled16steps-5s-Diffusers** | 5 second Diffusion distilled to 16 steps | 6× faster inference, minimal quality loss |
| **ai-forever/Kandinsky-5.0-T2V-Lite-distilled16steps-10s-Diffusers** | 10 second Diffusion distilled to 16 steps | 6× faster inference, minimal quality loss |
| **ai-forever/Kandinsky-5.0-T2V-Lite-pretrain-5s-Diffusers** | 5 second Base pretrained model | Research and fine-tuning |
| **ai-forever/Kandinsky-5.0-T2V-Lite-pretrain-10s-Diffusers** | 10 second Base pretrained model | Research and fine-tuning |
All models are available in 5-second and 10-second video generation versions.
## Kandinsky5T2VPipeline
[[autodoc]] Kandinsky5T2VPipeline
- all
- __call__
## Usage Examples
### Basic Text-to-Video Generation
```python
import torch
from diffusers import Kandinsky5T2VPipeline
from diffusers.utils import export_to_video
# Load the pipeline
model_id = "ai-forever/Kandinsky-5.0-T2V-Lite-sft-5s-Diffusers"
pipe = Kandinsky5T2VPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
pipe = pipe.to("cuda")
# Generate video
prompt = "A cat and a dog baking a cake together in a kitchen."
negative_prompt = "Static, 2D cartoon, cartoon, 2d animation, paintings, images, worst quality, low quality, ugly, deformed, walking backwards"
output = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
height=512,
width=768,
num_frames=121, # ~5 seconds at 24fps
num_inference_steps=50,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "output.mp4", fps=24, quality=9)
```
### 10 second Models
**⚠️ Warning!** all 10 second models should be used with Flex attention and max-autotune-no-cudagraphs compilation:
```python
pipe = Kandinsky5T2VPipeline.from_pretrained(
"ai-forever/Kandinsky-5.0-T2V-Lite-sft-10s-Diffusers",
torch_dtype=torch.bfloat16
)
pipe = pipe.to("cuda")
pipe.transformer.set_attention_backend(
"flex"
) # <--- Set attention backend to Flex
pipe.transformer.compile(
mode="max-autotune-no-cudagraphs",
dynamic=True
) # <--- Compile with max-autotune-no-cudagraphs
prompt = "A cat and a dog baking a cake together in a kitchen."
negative_prompt = "Static, 2D cartoon, cartoon, 2d animation, paintings, images, worst quality, low quality, ugly, deformed, walking backwards"
output = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
height=512,
width=768,
num_frames=241,
num_inference_steps=50,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "output.mp4", fps=24, quality=9)
```
### Diffusion Distilled model
**⚠️ Warning!** all nocfg and diffusion distilled models should be inferred without CFG (```guidance_scale=1.0```):
```python
model_id = "ai-forever/Kandinsky-5.0-T2V-Lite-distilled16steps-5s-Diffusers"
pipe = Kandinsky5T2VPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
pipe = pipe.to("cuda")
output = pipe(
prompt="A beautiful sunset over mountains",
num_inference_steps=16, # <--- Model is distilled in 16 steps
guidance_scale=1.0, # <--- no CFG
).frames[0]
export_to_video(output, "output.mp4", fps=24, quality=9)
```
## Citation
```bibtex
@misc{kandinsky2025,
author = {Alexey Letunovskiy and Maria Kovaleva and Ivan Kirillov and Lev Novitskiy and Denis Koposov and
Dmitrii Mikhailov and Anna Averchenkova and Andrey Shutkin and Julia Agafonova and Olga Kim and
Anastasiia Kargapoltseva and Nikita Kiselev and Vladimir Arkhipkin and Vladimir Korviakov and
Nikolai Gerasimenko and Denis Parkhomenko and Anna Dmitrienko and Anastasia Maltseva and
Kirill Chernyshev and Ilia Vasiliev and Viacheslav Vasilev and Vladimir Polovnikov and
Yury Kolabushin and Alexander Belykh and Mikhail Mamaev and Anastasia Aliaskina and
Tatiana Nikulina and Polina Gavrilova and Denis Dimitrov},
title = {Kandinsky 5.0: A family of diffusion models for Video & Image generation},
howpublished = {\url{https://github.com/ai-forever/Kandinsky-5}},
year = 2025
}
```
+131
View File
@@ -0,0 +1,131 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# PRX
PRX generates high-quality images from text using a simplified MMDIT architecture where text tokens don't update through transformer blocks. It employs flow matching with discrete scheduling for efficient sampling and uses Google's T5Gemma-2B-2B-UL2 model for multi-language text encoding. The ~1.3B parameter transformer delivers fast inference without sacrificing quality. You can choose between Flux VAE (8x compression, 16 latent channels) for balanced quality and speed or DC-AE (32x compression, 32 latent channels) for latent compression and faster processing.
## Available models
PRX offers multiple variants with different VAE configurations, each optimized for specific resolutions. Base models excel with detailed prompts, capturing complex compositions and subtle details. Fine-tuned models trained on the [Alchemist dataset](https://huggingface.co/datasets/yandex/alchemist) improve aesthetic quality, especially with simpler prompts.
| Model | Resolution | Fine-tuned | Distilled | Description | Suggested prompts | Suggested parameters | Recommended dtype |
|:-----:|:-----------------:|:----------:|:----------:|:----------:|:----------:|:----------:|:----------:|
| [`Photoroom/prx-256-t2i`](https://huggingface.co/Photoroom/prx-256-t2i)| 256 | No | No | Base model pre-trained at 256 with Flux VAE|Works best with detailed prompts in natural language|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-256-t2i-sft`](https://huggingface.co/Photoroom/prx-256-t2i-sft)| 512 | Yes | No | Fine-tuned on the [Alchemist dataset](https://huggingface.co/datasets/yandex/alchemist) dataset with Flux VAE | Can handle less detailed prompts|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i`](https://huggingface.co/Photoroom/prx-512-t2i)| 512 | No | No | Base model pre-trained at 512 with Flux VAE |Works best with detailed prompts in natural language|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i-sft`](https://huggingface.co/Photoroom/prx-512-t2i-sft)| 512 | Yes | No | Fine-tuned on the [Alchemist dataset](https://huggingface.co/datasets/yandex/alchemist) dataset with Flux VAE | Can handle less detailed prompts in natural language|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i-sft-distilled`](https://huggingface.co/Photoroom/prx-512-t2i-sft-distilled)| 512 | Yes | Yes | 8-step distilled model from [`Photoroom/prx-512-t2i-sft`](https://huggingface.co/Photoroom/prx-512-t2i-sft) | Can handle less detailed prompts in natural language|8 steps, cfg=1.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i-dc-ae`](https://huggingface.co/Photoroom/prx-512-t2i-dc-ae)| 512 | No | No | Base model pre-trained at 512 with [Deep Compression Autoencoder (DC-AE)](https://hanlab.mit.edu/projects/dc-ae)|Works best with detailed prompts in natural language|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i-dc-ae-sft`](https://huggingface.co/Photoroom/prx-512-t2i-dc-ae-sft)| 512 | Yes | No | Fine-tuned on the [Alchemist dataset](https://huggingface.co/datasets/yandex/alchemist) dataset with [Deep Compression Autoencoder (DC-AE)](https://hanlab.mit.edu/projects/dc-ae) | Can handle less detailed prompts in natural language|28 steps, cfg=5.0| `torch.bfloat16` |
| [`Photoroom/prx-512-t2i-dc-ae-sft-distilled`](https://huggingface.co/Photoroom/prx-512-t2i-dc-ae-sft-distilled)| 512 | Yes | Yes | 8-step distilled model from [`Photoroom/prx-512-t2i-dc-ae-sft-distilled`](https://huggingface.co/Photoroom/prx-512-t2i-dc-ae-sft-distilled) | Can handle less detailed prompts in natural language|8 steps, cfg=1.0| `torch.bfloat16` |s
Refer to [this](https://huggingface.co/collections/Photoroom/prx-models-68e66254c202ebfab99ad38e) collection for more information.
## Loading the pipeline
Load the pipeline with [`~DiffusionPipeline.from_pretrained`].
```py
from diffusers.pipelines.prx import PRXPipeline
# Load pipeline - VAE and text encoder will be loaded from HuggingFace
pipe = PRXPipeline.from_pretrained("Photoroom/prx-512-t2i-sft", torch_dtype=torch.bfloat16)
pipe.to("cuda")
prompt = "A front-facing portrait of a lion the golden savanna at sunset."
image = pipe(prompt, num_inference_steps=28, guidance_scale=5.0).images[0]
image.save("prx_output.png")
```
### Manual Component Loading
Load components individually to customize the pipeline for instance to use quantized models.
```py
import torch
from diffusers.pipelines.prx import PRXPipeline
from diffusers.models import AutoencoderKL, AutoencoderDC
from diffusers.models.transformers.transformer_prx import PRXTransformer2DModel
from diffusers.schedulers import FlowMatchEulerDiscreteScheduler
from transformers import T5GemmaModel, GemmaTokenizerFast
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as BitsAndBytesConfig
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
# Load transformer
transformer = PRXTransformer2DModel.from_pretrained(
"checkpoints/prx-512-t2i-sft",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
)
# Load scheduler
scheduler = FlowMatchEulerDiscreteScheduler.from_pretrained(
"checkpoints/prx-512-t2i-sft", subfolder="scheduler"
)
# Load T5Gemma text encoder
t5gemma_model = T5GemmaModel.from_pretrained("google/t5gemma-2b-2b-ul2",
quantization_config=quant_config,
torch_dtype=torch.bfloat16)
text_encoder = t5gemma_model.encoder.to(dtype=torch.bfloat16)
tokenizer = GemmaTokenizerFast.from_pretrained("google/t5gemma-2b-2b-ul2")
tokenizer.model_max_length = 256
# Load VAE - choose either Flux VAE or DC-AE
# Flux VAE
vae = AutoencoderKL.from_pretrained("black-forest-labs/FLUX.1-dev",
subfolder="vae",
quantization_config=quant_config,
torch_dtype=torch.bfloat16)
pipe = PRXPipeline(
transformer=transformer,
scheduler=scheduler,
text_encoder=text_encoder,
tokenizer=tokenizer,
vae=vae
)
pipe.to("cuda")
```
## Memory Optimization
For memory-constrained environments:
```py
import torch
from diffusers.pipelines.prx import PRXPipeline
pipe = PRXPipeline.from_pretrained("Photoroom/prx-512-t2i-sft", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload() # Offload components to CPU when not in use
# Or use sequential CPU offload for even lower memory
pipe.enable_sequential_cpu_offload()
```
## PRXPipeline
[[autodoc]] PRXPipeline
- all
- __call__
## PRXPipelineOutput
[[autodoc]] pipelines.prx.pipeline_output.PRXPipelineOutput
@@ -21,7 +21,7 @@ The Stable Diffusion model can also infer depth based on an image using [MiDaS](
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionDepth2ImgPipeline
@@ -21,14 +21,14 @@ The Stable Diffusion model can also be applied to inpainting which lets you edit
## Tips
It is recommended to use this pipeline with checkpoints that have been specifically fine-tuned for inpainting, such
as [runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting). Default
as [stable-diffusion-v1-5/stable-diffusion-inpainting](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting). Default
text-to-image Stable Diffusion checkpoints, such as
[stable-diffusion-v1-5/stable-diffusion-v1-5](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) are also compatible but they might be less performant.
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionInpaintPipeline
@@ -17,7 +17,7 @@ The Stable Diffusion latent upscaler model was created by [Katherine Crowson](ht
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionLatentUpscalePipeline
@@ -22,7 +22,7 @@ Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B data
For more details about how Stable Diffusion works and how it differs from the base latent diffusion model, take a look at the Stability AI [announcement](https://stability.ai/blog/stable-diffusion-announcement) and our own [blog post](https://huggingface.co/blog/stable_diffusion#how-does-stable-diffusion-work) for more technical details.
You can find the original codebase for Stable Diffusion v1.0 at [CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion) and Stable Diffusion v2.0 at [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion) as well as their original scripts for various tasks. Additional official checkpoints for the different Stable Diffusion versions and tasks can be found on the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations. Explore these organizations to find the best checkpoint for your use-case!
You can find the original codebase for Stable Diffusion v1.0 at [CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion) and Stable Diffusion v2.0 at [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion) as well as their original scripts for various tasks. Additional official checkpoints for the different Stable Diffusion versions and tasks can be found on the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations. Explore these organizations to find the best checkpoint for your use-case!
The table below summarizes the available Stable Diffusion pipelines, their supported tasks, and an interactive demo:
@@ -64,7 +64,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<a href="./inpaint">StableDiffusionInpaint</a>
</td>
<td class="px-4 py-2 text-gray-700">inpainting</td>
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/runwayml/stable-diffusion-inpainting"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/stable-diffusion-v1-5/stable-diffusion-inpainting"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
@@ -36,7 +36,7 @@ Here are some examples for how to use Stable Diffusion 2 for each task:
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## Text-to-image
@@ -25,7 +25,7 @@ The abstract from the paper is:
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionPipeline
@@ -21,7 +21,7 @@ The Stable Diffusion upscaler diffusion model was created by the researchers and
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionUpscalePipeline
@@ -21,6 +21,7 @@ Refer to the table below for an overview of the available attention families and
| attention family | main feature |
|---|---|
| FlashAttention | minimizes memory reads/writes through tiling and recomputation |
| AI Tensor Engine for ROCm | FlashAttention implementation optimized for AMD ROCm accelerators |
| SageAttention | quantizes attention to int8 |
| PyTorch native | built-in PyTorch implementation using [scaled_dot_product_attention](./fp16#scaled-dot-product-attention) |
| xFormers | memory-efficient attention with support for various attention kernels |
@@ -139,6 +140,7 @@ Refer to the table below for a complete list of available attention backends and
| `_native_xla` | [PyTorch native](https://docs.pytorch.org/docs/stable/generated/torch.nn.attention.SDPBackend.html#torch.nn.attention.SDPBackend) | XLA-optimized attention |
| `flash` | [FlashAttention](https://github.com/Dao-AILab/flash-attention) | FlashAttention-2 |
| `flash_varlen` | [FlashAttention](https://github.com/Dao-AILab/flash-attention) | Variable length FlashAttention |
| `aiter` | [AI Tensor Engine for ROCm](https://github.com/ROCm/aiter) | FlashAttention for AMD ROCm |
| `_flash_3` | [FlashAttention](https://github.com/Dao-AILab/flash-attention) | FlashAttention-3 |
| `_flash_varlen_3` | [FlashAttention](https://github.com/Dao-AILab/flash-attention) | Variable length FlashAttention-3 |
| `_flash_3_hub` | [FlashAttention](https://github.com/Dao-AILab/flash-attention) | FlashAttention-3 from kernels |
+2 -2
View File
@@ -16,12 +16,12 @@ pipeline.unet.config["in_channels"]
4
```
Inpainting requires 9 channels in the input sample. You can check this value in a pretrained inpainting model like [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting):
Inpainting requires 9 channels in the input sample. You can check this value in a pretrained inpainting model like [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting):
```py
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", use_safetensors=True)
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting", use_safetensors=True)
pipeline.unet.config["in_channels"]
9
```
@@ -215,7 +215,7 @@ from diffusers import AutoPipelineForInpainting, LCMScheduler
from diffusers.utils import load_image, make_image_grid
pipe = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
variant="fp16",
).to("cuda")
+15 -15
View File
@@ -112,7 +112,7 @@ blurred_mask
## Popular models
[Stable Diffusion Inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2 Inpainting](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
[Stable Diffusion Inpainting](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2 Inpainting](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
### Stable Diffusion Inpainting
@@ -124,7 +124,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -244,7 +244,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
```
</hfoption>
<hfoption id="runwayml/stable-diffusion-inpainting">
<hfoption id="stable-diffusion-v1-5/stable-diffusion-inpainting">
```py
import torch
@@ -252,7 +252,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -278,7 +278,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint-specific.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">runwayml/stable-diffusion-inpainting</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">stable-diffusion-v1-5/stable-diffusion-inpainting</figcaption>
</div>
</div>
@@ -308,7 +308,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
```
</hfoption>
<hfoption id="runwayml/stable-diffusion-inpaint">
<hfoption id="stable-diffusion-v1-5/stable-diffusion-inpaint">
```py
import torch
@@ -316,7 +316,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -340,7 +340,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/specific-inpaint-basic.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">runwayml/stable-diffusion-inpainting</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">stable-diffusion-v1-5/stable-diffusion-inpainting</figcaption>
</div>
</div>
@@ -358,7 +358,7 @@ from diffusers.utils import load_image, make_image_grid
device = "cuda"
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
variant="fp16"
)
@@ -396,7 +396,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -441,7 +441,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -481,7 +481,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -606,7 +606,7 @@ from diffusers import AutoPipelineForInpainting, AutoPipelineForImage2Image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -683,7 +683,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16,
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16,
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -714,7 +714,7 @@ controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_inpai
# pass ControlNet to the pipeline
pipeline = StableDiffusionControlNetInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
+1 -1
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@@ -173,7 +173,7 @@ mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
path = "runwayml/stable-diffusion-inpainting"
path = "stable-diffusion-v1-5/stable-diffusion-inpainting"
run_compile = True # Set True / False
+2 -2
View File
@@ -28,12 +28,12 @@ pipeline.unet.config["in_channels"]
4
```
인페인팅은 입력 샘플에 9개의 채널이 필요합니다. [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting)와 같은 사전학습된 인페인팅 모델에서 이 값을 확인할 수 있습니다:
인페인팅은 입력 샘플에 9개의 채널이 필요합니다. [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting)와 같은 사전학습된 인페인팅 모델에서 이 값을 확인할 수 있습니다:
```py
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting")
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting")
pipeline.unet.config["in_channels"]
9
```
+2 -11
View File
@@ -14,7 +14,7 @@ specific language governing permissions and limitations under the License.
[[open-in-colab]]
[`StableDiffusionInpaintPipeline`]은 마스크와 텍스트 프롬프트를 제공하여 이미지의 특정 부분을 편집할 수 있도록 합니다. 이 기능은 인페인팅 작업을 위해 특별히 훈련된 [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting)과 같은 Stable Diffusion 버전을 사용합니다.
[`StableDiffusionInpaintPipeline`]은 마스크와 텍스트 프롬프트를 제공하여 이미지의 특정 부분을 편집할 수 있도록 합니다. 이 기능은 인페인팅 작업을 위해 특별히 훈련된 [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting)과 같은 Stable Diffusion 버전을 사용합니다.
먼저 [`StableDiffusionInpaintPipeline`] 인스턴스를 불러옵니다:
@@ -27,7 +27,7 @@ from io import BytesIO
from diffusers import StableDiffusionInpaintPipeline
pipeline = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
)
pipeline = pipeline.to("cuda")
@@ -61,12 +61,3 @@ image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
> [!WARNING]
> 이전의 실험적인 인페인팅 구현에서는 품질이 낮은 다른 프로세스를 사용했습니다. 이전 버전과의 호환성을 보장하기 위해 새 모델이 포함되지 않은 사전학습된 파이프라인을 불러오면 이전 인페인팅 방법이 계속 적용됩니다.
아래 Space에서 이미지 인페인팅을 직접 해보세요!
<iframe
src="https://runwayml-stable-diffusion-inpainting.hf.space"
frameborder="0"
width="850"
height="500"
></iframe>
+2 -2
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@@ -16,12 +16,12 @@ pipeline.unet.config["in_channels"]
4
```
而图像修复任务需要输入样本具有9个通道。您可以在 [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting) 这样的预训练修复模型中验证此参数:
而图像修复任务需要输入样本具有9个通道。您可以在 [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting) 这样的预训练修复模型中验证此参数:
```python
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", use_safetensors=True)
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting", use_safetensors=True)
pipeline.unet.config["in_channels"]
9
```
+1 -1
View File
@@ -1328,7 +1328,7 @@ model = CLIPSegForImageSegmentation.from_pretrained("CIDAS/clipseg-rd64-refined"
# Load Stable Diffusion Inpainting Pipeline with custom pipeline
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
custom_pipeline="text_inpainting",
segmentation_model=model,
segmentation_processor=processor
@@ -126,7 +126,7 @@ EXAMPLE_DOC_STRING = """
... "lllyasviel/control_v11p_sd15_inpaint", torch_dtype=torch.float16
... )
>>> pipe = StableDiffusionControlNetInpaintPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config)
@@ -347,7 +347,7 @@ class AdaptiveMaskInpaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -429,8 +429,8 @@ class AdaptiveMaskInpaintPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -970,7 +970,7 @@ class AdaptiveMaskInpaintPipeline(
>>> default_mask_image = download_image(mask_url).resize((512, 512))
>>> pipe = AdaptiveMaskInpaintPipeline.from_pretrained(
... "runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16
... )
>>> pipe = pipe.to("cuda")
@@ -1095,7 +1095,7 @@ class AdaptiveMaskInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
@@ -62,7 +62,7 @@ class ComposableStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin)
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -145,8 +145,8 @@ class ComposableStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin)
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -1276,7 +1276,7 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
+1 -1
View File
@@ -678,7 +678,7 @@ class StableDiffusionHDPainterPipeline(StableDiffusionInpaintPipeline):
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+1 -1
View File
@@ -78,7 +78,7 @@ class ImageToImageInpaintingPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+3 -3
View File
@@ -86,7 +86,7 @@ class InstaFlowPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -165,8 +165,8 @@ class InstaFlowPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+3 -3
View File
@@ -166,7 +166,7 @@ class IPAdapterFaceIDStableDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -247,8 +247,8 @@ class IPAdapterFaceIDStableDiffusionPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -414,7 +414,7 @@ class StableDiffusionHighResFixPipeline(StableDiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -222,7 +222,7 @@ class LatentConsistencyModelWalkPipeline(
supports [`LCMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
+3 -3
View File
@@ -302,7 +302,7 @@ class LLMGroundedDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -392,8 +392,8 @@ class LLMGroundedDiffusionPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+2 -2
View File
@@ -552,8 +552,8 @@ class StableDiffusionLongPromptWeightingPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1765,7 +1765,7 @@ class SDXLLongPromptWeightingPipeline(
# Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != num_channels_unet:
+2 -2
View File
@@ -3729,8 +3729,8 @@ class MatryoshkaPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -78,7 +78,7 @@ class MultilingualStableDiffusion(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -1607,7 +1607,7 @@ class KolorsControlNetInpaintPipeline(
# 9. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+3 -3
View File
@@ -135,7 +135,7 @@ class FabricPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
"""
@@ -163,8 +163,8 @@ class FabricPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1487,7 +1487,7 @@ class KolorsInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+3 -3
View File
@@ -106,7 +106,7 @@ class Prompt2PromptPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -187,8 +187,8 @@ class Prompt2PromptPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1730,7 +1730,7 @@ class StyleAlignedSDXLPipeline(
# Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != num_channels_unet:
@@ -59,7 +59,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
@@ -392,7 +392,7 @@ class StableDiffusionBoxDiffPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -473,8 +473,8 @@ class StableDiffusionBoxDiffPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -42,7 +42,7 @@ EXAMPLE_DOC_STRING = """
```py
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
>>> image = pipe(prompt).images[0]
@@ -359,7 +359,7 @@ class StableDiffusionPAGPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -440,8 +440,8 @@ class StableDiffusionPAGPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -100,7 +100,7 @@ class StableDiffusionUpscaleLDM3DPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -2042,7 +2042,7 @@ class StableDiffusionXL_AE_Pipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
@@ -188,7 +188,7 @@ class StableDiffusionXLControlNetAdapterPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -330,7 +330,7 @@ class StableDiffusionXLControlNetAdapterInpaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
requires_aesthetics_score (`bool`, *optional*, defaults to `"False"`):
@@ -1569,7 +1569,7 @@ class StableDiffusionXLControlNetAdapterInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+4 -4
View File
@@ -46,7 +46,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
@@ -86,7 +86,7 @@ class Zero1to3StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
cc_projection ([`CCProjection`]):
@@ -164,8 +164,8 @@ class Zero1to3StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -288,7 +288,7 @@ class RerenderAVideoPipeline(StableDiffusionControlNetImg2ImgPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+1 -1
View File
@@ -54,7 +54,7 @@ EXAMPLE_DOC_STRING = """
>>> # load control net and stable diffusion v1-5
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> # speed up diffusion process with faster scheduler and memory optimization
@@ -158,7 +158,7 @@ EXAMPLE_DOC_STRING = """
>>> # load control net and stable diffusion v1-5
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> # speed up diffusion process with faster scheduler and memory optimization
@@ -64,7 +64,7 @@ class StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+1 -1
View File
@@ -114,7 +114,7 @@ class SdeDragPipeline(DiffusionPipeline):
>>> from diffusers import DDIMScheduler, DiffusionPipeline
>>> # Load the pipeline
>>> model_path = "runwayml/stable-diffusion-v1-5"
>>> model_path = "stable-diffusion-v1-5/stable-diffusion-v1-5"
>>> scheduler = DDIMScheduler.from_pretrained(model_path, subfolder="scheduler")
>>> pipe = DiffusionPipeline.from_pretrained(model_path, scheduler=scheduler, custom_pipeline="sde_drag")
>>> pipe.to('cuda')
@@ -46,7 +46,7 @@ class StableDiffusionComparisonPipeline(DiffusionPipeline, StableDiffusionMixin)
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionMegaSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -36,7 +36,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe_controlnet = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
controlnet=controlnet,
safety_checker=None,
torch_dtype=torch.float16
@@ -81,7 +81,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-seg", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
)
>>> pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
@@ -80,7 +80,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-seg", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetInpaintImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
)
>>> pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
@@ -37,7 +37,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetReferencePipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
controlnet=controlnet,
safety_checker=None,
torch_dtype=torch.float16
+4 -4
View File
@@ -43,7 +43,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_ipex")
>>> pipe = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_ipex")
>>> # For Float32
>>> pipe.prepare_for_ipex(prompt, dtype=torch.float32, height=512, width=512) #value of image height/width should be consistent with the pipeline inference
@@ -85,7 +85,7 @@ class StableDiffusionIPEXPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -161,8 +161,8 @@ class StableDiffusionIPEXPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -47,7 +47,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionMegaSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -46,7 +46,7 @@ EXAMPLE_DOC_STRING = """
>>> input_image = load_image("https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png")
>>> pipe = StableDiffusionReferencePipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
safety_checker=None,
torch_dtype=torch.float16
).to('cuda:0')
@@ -112,7 +112,7 @@ class StableDiffusionReferencePipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -194,8 +194,8 @@ class StableDiffusionReferencePipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -167,7 +167,7 @@ class StableDiffusionRepaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -249,8 +249,8 @@ class StableDiffusionRepaintPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -678,7 +678,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -766,8 +766,8 @@ class TensorRTStableDiffusionImg2ImgPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -682,7 +682,7 @@ class TensorRTStableDiffusionInpaintPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -770,8 +770,8 @@ class TensorRTStableDiffusionInpaintPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -594,7 +594,7 @@ class TensorRTStableDiffusionPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -682,8 +682,8 @@ class TensorRTStableDiffusionPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -52,7 +52,7 @@ class TextInpainting(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -1223,7 +1223,7 @@ class AnyTextPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -5,7 +5,7 @@ This script was added by @thedarkzeno .
Please note that this script is not actively maintained, you can open an issue and tag @thedarkzeno or @patil-suraj though.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export OUTPUT_DIR="path-to-save-model"
@@ -29,7 +29,7 @@ Prior-preservation is used to avoid overfitting and language-drift. Refer to the
According to the paper, it's recommended to generate `num_epochs * num_samples` images for prior-preservation. 200-300 works well for most cases.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -60,7 +60,7 @@ With the help of gradient checkpointing and the 8-bit optimizer from bitsandbyte
To install `bitandbytes` please refer to this [readme](https://github.com/TimDettmers/bitsandbytes#requirements--installation).
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -92,7 +92,7 @@ Pass the `--train_text_encoder` argument to the script to enable training `text_
___Note: Training text encoder requires more memory, with this option the training won't fit on 16GB GPU. It needs at least 24GB VRAM.___
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -55,7 +55,7 @@ The Accelerate launch command is used to train a model using multiple GPUs and m
```
accelerate launch --mixed_precision "fp16" \
tutorial_train_ip-adapter.py \
--pretrained_model_name_or_path="runwayml/stable-diffusion-v1-5/" \
--pretrained_model_name_or_path="stable-diffusion-v1-5/stable-diffusion-v1-5/" \
--image_encoder_path="{image_encoder_path}" \
--data_json_file="{data.json}" \
--data_root_path="{image_path}" \
@@ -73,7 +73,7 @@ tutorial_train_ip-adapter.py \
```
accelerate launch --num_processes 8 --multi_gpu --mixed_precision "fp16" \
tutorial_train_ip-adapter.py \
--pretrained_model_name_or_path="runwayml/stable-diffusion-v1-5/" \
--pretrained_model_name_or_path="stable-diffusion-v1-5/stable-diffusion-v1-5/" \
--image_encoder_path="{image_encoder_path}" \
--data_json_file="{data.json}" \
--data_root_path="{image_path}" \
@@ -27,7 +27,7 @@ You can build multiple datasets for every subject and upload them to the 🤗 hu
Before launching the training script, make sure to select the inpainting the target model, the output directory and the 🤗 datasets.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export OUTPUT_DIR="path-to-save-model"
export DATASET_1="gzguevara/mr_potato_head_masked"
@@ -177,7 +177,7 @@ class PromptDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -238,7 +238,7 @@ def parse_args() -> argparse.Namespace:
# EXAMPLE USAGE:
#
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "runwayml/stable-diffusion-v1-5" --subfolder "vae" --input_image "foo.png"
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "stable-diffusion-v1-5/stable-diffusion-v1-5" --subfolder "vae" --input_image "foo.png"
#
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "madebyollin/taesd" --use_tiny_nn --input_image "foo.png"
#
+7 -4
View File
@@ -24,7 +24,8 @@ args = args.parse_args()
def _extract_into_tensor(arr, timesteps, broadcast_shape):
# from: https://github.com/openai/guided-diffusion/blob/22e0df8183507e13a7813f8d38d51b072ca1e67c/guided_diffusion/gaussian_diffusion.py#L895 """
# from: https://github.com/openai/guided-diffusion/blob/22e0df8183507e13a7813f8d38d51b072ca1e67c/guided_diffusion/gaussian_diffusion.py#L895
# """
res = arr[timesteps].float()
dims_to_append = len(broadcast_shape) - len(res.shape)
return res[(...,) + (None,) * dims_to_append]
@@ -507,7 +508,9 @@ def rename_state_dict(sd, embedding):
# encode with stable diffusion vae
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16
)
pipe.vae.cuda()
# construct original decoder with jitted model
@@ -1090,7 +1093,7 @@ def new_constructor(self, **kwargs):
Encoder.__init__ = new_constructor
vae = AutoencoderKL.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="vae")
vae = AutoencoderKL.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="vae")
consistency_vae = ConsistencyDecoderVAE(
encoder_args=vae.encoder.constructor_arguments,
decoder_args=unet.config,
@@ -1117,7 +1120,7 @@ print((sample_consistency_orig - sample_consistency_new_3).abs().sum())
print("running with diffusers pipeline")
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", vae=consistency_vae, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-v1-5", vae=consistency_vae, torch_dtype=torch.float16
)
pipe.to("cuda")
File diff suppressed because it is too large Load Diff
+345
View File
@@ -0,0 +1,345 @@
#!/usr/bin/env python3
"""
Script to convert PRX checkpoint from original codebase to diffusers format.
"""
import argparse
import json
import os
import sys
from dataclasses import asdict, dataclass
from typing import Dict, Tuple
import torch
from safetensors.torch import save_file
from diffusers.models.transformers.transformer_prx import PRXTransformer2DModel
from diffusers.pipelines.prx import PRXPipeline
DEFAULT_RESOLUTION = 512
@dataclass(frozen=True)
class PRXBase:
context_in_dim: int = 2304
hidden_size: int = 1792
mlp_ratio: float = 3.5
num_heads: int = 28
depth: int = 16
axes_dim: Tuple[int, int] = (32, 32)
theta: int = 10_000
time_factor: float = 1000.0
time_max_period: int = 10_000
@dataclass(frozen=True)
class PRXFlux(PRXBase):
in_channels: int = 16
patch_size: int = 2
@dataclass(frozen=True)
class PRXDCAE(PRXBase):
in_channels: int = 32
patch_size: int = 1
def build_config(vae_type: str) -> Tuple[dict, int]:
if vae_type == "flux":
cfg = PRXFlux()
elif vae_type == "dc-ae":
cfg = PRXDCAE()
else:
raise ValueError(f"Unsupported VAE type: {vae_type}. Use 'flux' or 'dc-ae'")
config_dict = asdict(cfg)
config_dict["axes_dim"] = list(config_dict["axes_dim"]) # type: ignore[index]
return config_dict
def create_parameter_mapping(depth: int) -> dict:
"""Create mapping from old parameter names to new diffusers names."""
# Key mappings for structural changes
mapping = {}
# Map old structure (layers in PRXBlock) to new structure (layers in PRXAttention)
for i in range(depth):
# QKV projections moved to attention module
mapping[f"blocks.{i}.img_qkv_proj.weight"] = f"blocks.{i}.attention.img_qkv_proj.weight"
mapping[f"blocks.{i}.txt_kv_proj.weight"] = f"blocks.{i}.attention.txt_kv_proj.weight"
# QK norm moved to attention module and renamed to match Attention's qk_norm structure
mapping[f"blocks.{i}.qk_norm.query_norm.scale"] = f"blocks.{i}.attention.norm_q.weight"
mapping[f"blocks.{i}.qk_norm.key_norm.scale"] = f"blocks.{i}.attention.norm_k.weight"
mapping[f"blocks.{i}.qk_norm.query_norm.weight"] = f"blocks.{i}.attention.norm_q.weight"
mapping[f"blocks.{i}.qk_norm.key_norm.weight"] = f"blocks.{i}.attention.norm_k.weight"
# K norm for text tokens moved to attention module
mapping[f"blocks.{i}.k_norm.scale"] = f"blocks.{i}.attention.norm_added_k.weight"
mapping[f"blocks.{i}.k_norm.weight"] = f"blocks.{i}.attention.norm_added_k.weight"
# Attention output projection
mapping[f"blocks.{i}.attn_out.weight"] = f"blocks.{i}.attention.to_out.0.weight"
return mapping
def convert_checkpoint_parameters(old_state_dict: Dict[str, torch.Tensor], depth: int) -> Dict[str, torch.Tensor]:
"""Convert old checkpoint parameters to new diffusers format."""
print("Converting checkpoint parameters...")
mapping = create_parameter_mapping(depth)
converted_state_dict = {}
for key, value in old_state_dict.items():
new_key = key
# Apply specific mappings if needed
if key in mapping:
new_key = mapping[key]
print(f" Mapped: {key} -> {new_key}")
converted_state_dict[new_key] = value
print(f"✓ Converted {len(converted_state_dict)} parameters")
return converted_state_dict
def create_transformer_from_checkpoint(checkpoint_path: str, config: dict) -> PRXTransformer2DModel:
"""Create and load PRXTransformer2DModel from old checkpoint."""
print(f"Loading checkpoint from: {checkpoint_path}")
# Load old checkpoint
if not os.path.exists(checkpoint_path):
raise FileNotFoundError(f"Checkpoint not found: {checkpoint_path}")
old_checkpoint = torch.load(checkpoint_path, map_location="cpu")
# Handle different checkpoint formats
if isinstance(old_checkpoint, dict):
if "model" in old_checkpoint:
state_dict = old_checkpoint["model"]
elif "state_dict" in old_checkpoint:
state_dict = old_checkpoint["state_dict"]
else:
state_dict = old_checkpoint
else:
state_dict = old_checkpoint
print(f"✓ Loaded checkpoint with {len(state_dict)} parameters")
# Convert parameter names if needed
model_depth = int(config.get("depth", 16))
converted_state_dict = convert_checkpoint_parameters(state_dict, depth=model_depth)
# Create transformer with config
print("Creating PRXTransformer2DModel...")
transformer = PRXTransformer2DModel(**config)
# Load state dict
print("Loading converted parameters...")
missing_keys, unexpected_keys = transformer.load_state_dict(converted_state_dict, strict=False)
if missing_keys:
print(f"⚠ Missing keys: {missing_keys}")
if unexpected_keys:
print(f"⚠ Unexpected keys: {unexpected_keys}")
if not missing_keys and not unexpected_keys:
print("✓ All parameters loaded successfully!")
return transformer
def create_scheduler_config(output_path: str, shift: float):
"""Create FlowMatchEulerDiscreteScheduler config."""
scheduler_config = {"_class_name": "FlowMatchEulerDiscreteScheduler", "num_train_timesteps": 1000, "shift": shift}
scheduler_path = os.path.join(output_path, "scheduler")
os.makedirs(scheduler_path, exist_ok=True)
with open(os.path.join(scheduler_path, "scheduler_config.json"), "w") as f:
json.dump(scheduler_config, f, indent=2)
print("✓ Created scheduler config")
def download_and_save_vae(vae_type: str, output_path: str):
"""Download and save VAE to local directory."""
from diffusers import AutoencoderDC, AutoencoderKL
vae_path = os.path.join(output_path, "vae")
os.makedirs(vae_path, exist_ok=True)
if vae_type == "flux":
print("Downloading FLUX VAE from black-forest-labs/FLUX.1-dev...")
vae = AutoencoderKL.from_pretrained("black-forest-labs/FLUX.1-dev", subfolder="vae")
else: # dc-ae
print("Downloading DC-AE VAE from mit-han-lab/dc-ae-f32c32-sana-1.1-diffusers...")
vae = AutoencoderDC.from_pretrained("mit-han-lab/dc-ae-f32c32-sana-1.1-diffusers")
vae.save_pretrained(vae_path)
print(f"✓ Saved VAE to {vae_path}")
def download_and_save_text_encoder(output_path: str):
"""Download and save T5Gemma text encoder and tokenizer."""
from transformers import GemmaTokenizerFast
from transformers.models.t5gemma.modeling_t5gemma import T5GemmaModel
text_encoder_path = os.path.join(output_path, "text_encoder")
tokenizer_path = os.path.join(output_path, "tokenizer")
os.makedirs(text_encoder_path, exist_ok=True)
os.makedirs(tokenizer_path, exist_ok=True)
print("Downloading T5Gemma model from google/t5gemma-2b-2b-ul2...")
t5gemma_model = T5GemmaModel.from_pretrained("google/t5gemma-2b-2b-ul2")
# Extract and save only the encoder
t5gemma_encoder = t5gemma_model.encoder
t5gemma_encoder.save_pretrained(text_encoder_path)
print(f"✓ Saved T5GemmaEncoder to {text_encoder_path}")
print("Downloading tokenizer from google/t5gemma-2b-2b-ul2...")
tokenizer = GemmaTokenizerFast.from_pretrained("google/t5gemma-2b-2b-ul2")
tokenizer.model_max_length = 256
tokenizer.save_pretrained(tokenizer_path)
print(f"✓ Saved tokenizer to {tokenizer_path}")
def create_model_index(vae_type: str, default_image_size: int, output_path: str):
"""Create model_index.json for the pipeline."""
if vae_type == "flux":
vae_class = "AutoencoderKL"
else: # dc-ae
vae_class = "AutoencoderDC"
model_index = {
"_class_name": "PRXPipeline",
"_diffusers_version": "0.31.0.dev0",
"_name_or_path": os.path.basename(output_path),
"default_sample_size": default_image_size,
"scheduler": ["diffusers", "FlowMatchEulerDiscreteScheduler"],
"text_encoder": ["prx", "T5GemmaEncoder"],
"tokenizer": ["transformers", "GemmaTokenizerFast"],
"transformer": ["diffusers", "PRXTransformer2DModel"],
"vae": ["diffusers", vae_class],
}
model_index_path = os.path.join(output_path, "model_index.json")
with open(model_index_path, "w") as f:
json.dump(model_index, f, indent=2)
def main(args):
# Validate inputs
if not os.path.exists(args.checkpoint_path):
raise FileNotFoundError(f"Checkpoint not found: {args.checkpoint_path}")
config = build_config(args.vae_type)
# Create output directory
os.makedirs(args.output_path, exist_ok=True)
print(f"✓ Output directory: {args.output_path}")
# Create transformer from checkpoint
transformer = create_transformer_from_checkpoint(args.checkpoint_path, config)
# Save transformer
transformer_path = os.path.join(args.output_path, "transformer")
os.makedirs(transformer_path, exist_ok=True)
# Save config
with open(os.path.join(transformer_path, "config.json"), "w") as f:
json.dump(config, f, indent=2)
# Save model weights as safetensors
state_dict = transformer.state_dict()
save_file(state_dict, os.path.join(transformer_path, "diffusion_pytorch_model.safetensors"))
print(f"✓ Saved transformer to {transformer_path}")
# Create scheduler config
create_scheduler_config(args.output_path, args.shift)
download_and_save_vae(args.vae_type, args.output_path)
download_and_save_text_encoder(args.output_path)
# Create model_index.json
create_model_index(args.vae_type, args.resolution, args.output_path)
# Verify the pipeline can be loaded
try:
pipeline = PRXPipeline.from_pretrained(args.output_path)
print("Pipeline loaded successfully!")
print(f"Transformer: {type(pipeline.transformer).__name__}")
print(f"VAE: {type(pipeline.vae).__name__}")
print(f"Text Encoder: {type(pipeline.text_encoder).__name__}")
print(f"Scheduler: {type(pipeline.scheduler).__name__}")
# Display model info
num_params = sum(p.numel() for p in pipeline.transformer.parameters())
print(f"✓ Transformer parameters: {num_params:,}")
except Exception as e:
print(f"Pipeline verification failed: {e}")
return False
print("Conversion completed successfully!")
print(f"Converted pipeline saved to: {args.output_path}")
print(f"VAE type: {args.vae_type}")
return True
if __name__ == "__main__":
parser = argparse.ArgumentParser(description="Convert PRX checkpoint to diffusers format")
parser.add_argument(
"--checkpoint_path", type=str, required=True, help="Path to the original PRX checkpoint (.pth file )"
)
parser.add_argument(
"--output_path", type=str, required=True, help="Output directory for the converted diffusers pipeline"
)
parser.add_argument(
"--vae_type",
type=str,
choices=["flux", "dc-ae"],
required=True,
help="VAE type to use: 'flux' for AutoencoderKL (16 channels) or 'dc-ae' for AutoencoderDC (32 channels)",
)
parser.add_argument(
"--resolution",
type=int,
choices=[256, 512, 1024],
default=DEFAULT_RESOLUTION,
help="Target resolution for the model (256, 512, or 1024). Affects the transformer's sample_size.",
)
parser.add_argument(
"--shift",
type=float,
default=3.0,
help="Shift for the scheduler",
)
args = parser.parse_args()
try:
success = main(args)
if not success:
sys.exit(1)
except Exception as e:
print(f"Conversion failed: {e}")
import traceback
traceback.print_exc()
sys.exit(1)
+26
View File
@@ -149,7 +149,9 @@ else:
_import_structure["guiders"].extend(
[
"AdaptiveProjectedGuidance",
"AdaptiveProjectedMixGuidance",
"AutoGuidance",
"BaseGuidance",
"ClassifierFreeGuidance",
"ClassifierFreeZeroStarGuidance",
"FrequencyDecoupledGuidance",
@@ -184,6 +186,8 @@ else:
"AutoencoderKLAllegro",
"AutoencoderKLCogVideoX",
"AutoencoderKLCosmos",
"AutoencoderKLHunyuanImage",
"AutoencoderKLHunyuanImageRefiner",
"AutoencoderKLHunyuanVideo",
"AutoencoderKLLTXVideo",
"AutoencoderKLMagvit",
@@ -194,6 +198,7 @@ else:
"AutoencoderOobleck",
"AutoencoderTiny",
"AutoModel",
"BriaFiboTransformer2DModel",
"BriaTransformer2DModel",
"CacheMixin",
"ChromaTransformer2DModel",
@@ -216,10 +221,12 @@ else:
"HunyuanDiT2DControlNetModel",
"HunyuanDiT2DModel",
"HunyuanDiT2DMultiControlNetModel",
"HunyuanImageTransformer2DModel",
"HunyuanVideoFramepackTransformer3DModel",
"HunyuanVideoTransformer3DModel",
"I2VGenXLUNet",
"Kandinsky3UNet",
"Kandinsky5Transformer3DModel",
"LatteTransformer3DModel",
"LTXVideoTransformer3DModel",
"Lumina2Transformer2DModel",
@@ -233,6 +240,7 @@ else:
"ParallelConfig",
"PixArtTransformer2DModel",
"PriorTransformer",
"PRXTransformer2DModel",
"QwenImageControlNetModel",
"QwenImageMultiControlNetModel",
"QwenImageTransformer2DModel",
@@ -423,6 +431,7 @@ else:
"AuraFlowPipeline",
"BlipDiffusionControlNetPipeline",
"BlipDiffusionPipeline",
"BriaFiboPipeline",
"BriaPipeline",
"ChromaImg2ImgPipeline",
"ChromaPipeline",
@@ -460,6 +469,8 @@ else:
"HunyuanDiTControlNetPipeline",
"HunyuanDiTPAGPipeline",
"HunyuanDiTPipeline",
"HunyuanImagePipeline",
"HunyuanImageRefinerPipeline",
"HunyuanSkyreelsImageToVideoPipeline",
"HunyuanVideoFramepackPipeline",
"HunyuanVideoImageToVideoPipeline",
@@ -474,6 +485,7 @@ else:
"ImageTextPipelineOutput",
"Kandinsky3Img2ImgPipeline",
"Kandinsky3Pipeline",
"Kandinsky5T2VPipeline",
"KandinskyCombinedPipeline",
"KandinskyImg2ImgCombinedPipeline",
"KandinskyImg2ImgPipeline",
@@ -517,6 +529,7 @@ else:
"PixArtAlphaPipeline",
"PixArtSigmaPAGPipeline",
"PixArtSigmaPipeline",
"PRXPipeline",
"QwenImageControlNetInpaintPipeline",
"QwenImageControlNetPipeline",
"QwenImageEditInpaintPipeline",
@@ -845,7 +858,9 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
else:
from .guiders import (
AdaptiveProjectedGuidance,
AdaptiveProjectedMixGuidance,
AutoGuidance,
BaseGuidance,
ClassifierFreeGuidance,
ClassifierFreeZeroStarGuidance,
FrequencyDecoupledGuidance,
@@ -876,6 +891,8 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
AutoencoderKLAllegro,
AutoencoderKLCogVideoX,
AutoencoderKLCosmos,
AutoencoderKLHunyuanImage,
AutoencoderKLHunyuanImageRefiner,
AutoencoderKLHunyuanVideo,
AutoencoderKLLTXVideo,
AutoencoderKLMagvit,
@@ -886,6 +903,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
AutoencoderOobleck,
AutoencoderTiny,
AutoModel,
BriaFiboTransformer2DModel,
BriaTransformer2DModel,
CacheMixin,
ChromaTransformer2DModel,
@@ -908,10 +926,12 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
HunyuanDiT2DControlNetModel,
HunyuanDiT2DModel,
HunyuanDiT2DMultiControlNetModel,
HunyuanImageTransformer2DModel,
HunyuanVideoFramepackTransformer3DModel,
HunyuanVideoTransformer3DModel,
I2VGenXLUNet,
Kandinsky3UNet,
Kandinsky5Transformer3DModel,
LatteTransformer3DModel,
LTXVideoTransformer3DModel,
Lumina2Transformer2DModel,
@@ -925,6 +945,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
ParallelConfig,
PixArtTransformer2DModel,
PriorTransformer,
PRXTransformer2DModel,
QwenImageControlNetModel,
QwenImageMultiControlNetModel,
QwenImageTransformer2DModel,
@@ -1085,6 +1106,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
AudioLDM2UNet2DConditionModel,
AudioLDMPipeline,
AuraFlowPipeline,
BriaFiboPipeline,
BriaPipeline,
ChromaImg2ImgPipeline,
ChromaPipeline,
@@ -1122,6 +1144,8 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
HunyuanDiTControlNetPipeline,
HunyuanDiTPAGPipeline,
HunyuanDiTPipeline,
HunyuanImagePipeline,
HunyuanImageRefinerPipeline,
HunyuanSkyreelsImageToVideoPipeline,
HunyuanVideoFramepackPipeline,
HunyuanVideoImageToVideoPipeline,
@@ -1136,6 +1160,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
ImageTextPipelineOutput,
Kandinsky3Img2ImgPipeline,
Kandinsky3Pipeline,
Kandinsky5T2VPipeline,
KandinskyCombinedPipeline,
KandinskyImg2ImgCombinedPipeline,
KandinskyImg2ImgPipeline,
@@ -1179,6 +1204,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
PixArtAlphaPipeline,
PixArtSigmaPAGPipeline,
PixArtSigmaPipeline,
PRXPipeline,
QwenImageControlNetInpaintPipeline,
QwenImageControlNetPipeline,
QwenImageEditInpaintPipeline,

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