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8 Commits

Author SHA1 Message Date
DN6 39115294b3 update 2025-10-21 11:51:18 +05:30
Steven Liu 5b5fa49a89 [docs] Organize toctree by modality (#12514)
* reorganize

* fix

---------

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>
2025-10-21 10:18:54 +05:30
Fei Xie decfa3c9e1 Fix: Use incorrect temporary variable key when replacing adapter name… (#12502)
Fix: Use incorrect temporary variable key when replacing adapter name in state dict within load_lora_adapter function

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2025-10-20 15:45:37 -10:00
Dhruv Nair 48305755bf Raise warning instead of error when imports are missing for custom code (#12513)
update
2025-10-20 07:02:23 -10:00
dg845 7853bfbed7 Remove Qwen Image Redundant RoPE Cache (#12452)
Refactor QwenEmbedRope to only use the LRU cache for RoPE caching
2025-10-19 18:41:58 -07:00
Lev Novitskiy 23ebbb4bc8 Kandinsky 5 is finally in Diffusers! (#12478)
* add kandinsky5 transformer pipeline first version

---------

Co-authored-by: Álvaro Somoza <asomoza@users.noreply.github.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Charles <charles@huggingface.co>
2025-10-17 18:34:30 -10:00
Ali Imran 1b456bd5d5 docs: cleanup of runway model (#12503)
* cleanup of runway model

* quality fixes
2025-10-17 14:10:50 -07:00
Sayak Paul af769881d3 [tests] introduce VAETesterMixin to consolidate tests for slicing and tiling (#12374)
* up

* up

* up

* up

* up

* u[

* up

* up

* up
2025-10-17 12:02:29 +05:30
115 changed files with 2476 additions and 514 deletions
+1 -1
View File
@@ -171,7 +171,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
<tr style="border-top: 2px solid black">
<td>Text-guided Image Inpainting</td>
<td><a href="https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/inpaint">Stable Diffusion Inpainting</a></td>
<td><a href="https://huggingface.co/runwayml/stable-diffusion-inpainting"> runwayml/stable-diffusion-inpainting </a></td>
<td><a href="https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting"> stable-diffusion-v1-5/stable-diffusion-inpainting </a></td>
</tr>
<tr style="border-top: 2px solid black">
<td>Image Variation</td>
+258 -263
View File
@@ -1,5 +1,4 @@
- title: Get started
sections:
- sections:
- local: index
title: Diffusers
- local: installation
@@ -8,9 +7,8 @@
title: Quickstart
- local: stable_diffusion
title: Basic performance
- title: Pipelines
isExpanded: false
title: Get started
- isExpanded: false
sections:
- local: using-diffusers/loading
title: DiffusionPipeline
@@ -28,9 +26,8 @@
title: Model formats
- local: using-diffusers/push_to_hub
title: Sharing pipelines and models
- title: Adapters
isExpanded: false
title: Pipelines
- isExpanded: false
sections:
- local: tutorials/using_peft_for_inference
title: LoRA
@@ -44,9 +41,8 @@
title: DreamBooth
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- title: Inference
isExpanded: false
title: Adapters
- isExpanded: false
sections:
- local: using-diffusers/weighted_prompts
title: Prompting
@@ -56,9 +52,8 @@
title: Batch inference
- local: training/distributed_inference
title: Distributed inference
- title: Inference optimization
isExpanded: false
title: Inference
- isExpanded: false
sections:
- local: optimization/fp16
title: Accelerate inference
@@ -70,8 +65,7 @@
title: Reduce memory usage
- local: optimization/speed-memory-optims
title: Compiling and offloading quantized models
- title: Community optimizations
sections:
- sections:
- local: optimization/pruna
title: Pruna
- local: optimization/xformers
@@ -90,9 +84,9 @@
title: ParaAttention
- local: using-diffusers/image_quality
title: FreeU
- title: Hybrid Inference
isExpanded: false
title: Community optimizations
title: Inference optimization
- isExpanded: false
sections:
- local: hybrid_inference/overview
title: Overview
@@ -102,9 +96,8 @@
title: VAE Encode
- local: hybrid_inference/api_reference
title: API Reference
- title: Modular Diffusers
isExpanded: false
title: Hybrid Inference
- isExpanded: false
sections:
- local: modular_diffusers/overview
title: Overview
@@ -126,9 +119,8 @@
title: ComponentsManager
- local: modular_diffusers/guiders
title: Guiders
- title: Training
isExpanded: false
title: Modular Diffusers
- isExpanded: false
sections:
- local: training/overview
title: Overview
@@ -138,8 +130,7 @@
title: Adapt a model to a new task
- local: tutorials/basic_training
title: Train a diffusion model
- title: Models
sections:
- sections:
- local: training/unconditional_training
title: Unconditional image generation
- local: training/text2image
@@ -158,8 +149,8 @@
title: InstructPix2Pix
- local: training/cogvideox
title: CogVideoX
- title: Methods
sections:
title: Models
- sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
@@ -172,9 +163,9 @@
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
- title: Quantization
isExpanded: false
title: Methods
title: Training
- isExpanded: false
sections:
- local: quantization/overview
title: Getting started
@@ -188,9 +179,8 @@
title: quanto
- local: quantization/modelopt
title: NVIDIA ModelOpt
- title: Model accelerators and hardware
isExpanded: false
title: Quantization
- isExpanded: false
sections:
- local: optimization/onnx
title: ONNX
@@ -204,9 +194,8 @@
title: Intel Gaudi
- local: optimization/neuron
title: AWS Neuron
- title: Specific pipeline examples
isExpanded: false
title: Model accelerators and hardware
- isExpanded: false
sections:
- local: using-diffusers/consisid
title: ConsisID
@@ -232,12 +221,10 @@
title: Stable Video Diffusion
- local: using-diffusers/marigold_usage
title: Marigold Computer Vision
- title: Resources
isExpanded: false
title: Specific pipeline examples
- isExpanded: false
sections:
- title: Task recipes
sections:
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: using-diffusers/conditional_image_generation
@@ -252,6 +239,7 @@
title: Video generation
- local: using-diffusers/depth2img
title: Depth-to-image
title: Task recipes
- local: using-diffusers/write_own_pipeline
title: Understanding pipelines, models and schedulers
- local: community_projects
@@ -266,12 +254,10 @@
title: Diffusers' Ethical Guidelines
- local: conceptual/evaluation
title: Evaluating Diffusion Models
- title: API
isExpanded: false
title: Resources
- isExpanded: false
sections:
- title: Main Classes
sections:
- sections:
- local: api/configuration
title: Configuration
- local: api/logging
@@ -282,8 +268,8 @@
title: Quantization
- local: api/parallel
title: Parallel inference
- title: Modular
sections:
title: Main Classes
- sections:
- local: api/modular_diffusers/pipeline
title: Pipeline
- local: api/modular_diffusers/pipeline_blocks
@@ -294,8 +280,8 @@
title: Components and configs
- local: api/modular_diffusers/guiders
title: Guiders
- title: Loaders
sections:
title: Modular
- sections:
- local: api/loaders/ip_adapter
title: IP-Adapter
- local: api/loaders/lora
@@ -310,14 +296,13 @@
title: SD3Transformer2D
- local: api/loaders/peft
title: PEFT
- title: Models
sections:
title: Loaders
- sections:
- local: api/models/overview
title: Overview
- local: api/models/auto_model
title: AutoModel
- title: ControlNets
sections:
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_union
@@ -332,8 +317,8 @@
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- title: Transformers
sections:
title: ControlNets
- sections:
- local: api/models/allegro_transformer3d
title: AllegroTransformer3DModel
- local: api/models/aura_flow_transformer2d
@@ -396,8 +381,8 @@
title: TransformerTemporalModel
- local: api/models/wan_transformer_3d
title: WanTransformer3DModel
- title: UNets
sections:
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/unet
@@ -412,8 +397,8 @@
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
- title: VAEs
sections:
title: UNets
- sections:
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
@@ -446,210 +431,218 @@
title: Tiny AutoEncoder
- local: api/models/vq
title: VQModel
- title: Pipelines
sections:
title: VAEs
title: Models
- sections:
- local: api/pipelines/overview
title: Overview
- local: api/pipelines/allegro
title: Allegro
- local: api/pipelines/amused
title: aMUSEd
- local: api/pipelines/animatediff
title: AnimateDiff
- local: api/pipelines/attend_and_excite
title: Attend-and-Excite
- local: api/pipelines/audioldm
title: AudioLDM
- local: api/pipelines/audioldm2
title: AudioLDM 2
- local: api/pipelines/aura_flow
title: AuraFlow
- sections:
- local: api/pipelines/audioldm
title: AudioLDM
- local: api/pipelines/audioldm2
title: AudioLDM 2
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/stable_audio
title: Stable Audio
title: Audio
- local: api/pipelines/auto_pipeline
title: AutoPipeline
- local: api/pipelines/blip_diffusion
title: BLIP-Diffusion
- local: api/pipelines/bria_3_2
title: Bria 3.2
- local: api/pipelines/chroma
title: Chroma
- local: api/pipelines/cogvideox
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_flux
title: ControlNet with Flux.1
- local: api/pipelines/controlnet_hunyuandit
title: ControlNet with Hunyuan-DiT
- local: api/pipelines/controlnet_sd3
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/cosmos
title: Cosmos
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
title: DDIM
- local: api/pipelines/ddpm
title: DDPM
- local: api/pipelines/deepfloyd_if
title: DeepFloyd IF
- local: api/pipelines/diffedit
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/framepack
title: Framepack
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/hunyuan_video
title: HunyuanVideo
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/pix2pix
title: InstructPix2Pix
- local: api/pipelines/kandinsky
title: Kandinsky 2.1
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/kandinsky3
title: Kandinsky 3
- local: api/pipelines/kolors
title: Kolors
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/latte
title: Latte
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/mochi
title: Mochi
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
title: Paint by Example
- local: api/pipelines/pia
title: Personalized Image Animator (PIA)
- local: api/pipelines/pixart
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/qwenimage
title: QwenImage
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/skyreels_v2
title: SkyReels-V2
- local: api/pipelines/stable_audio
title: Stable Audio
- local: api/pipelines/stable_cascade
title: Stable Cascade
- title: Stable Diffusion
sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-image
- sections:
- local: api/pipelines/amused
title: aMUSEd
- local: api/pipelines/animatediff
title: AnimateDiff
- local: api/pipelines/attend_and_excite
title: Attend-and-Excite
- local: api/pipelines/aura_flow
title: AuraFlow
- local: api/pipelines/blip_diffusion
title: BLIP-Diffusion
- local: api/pipelines/bria_3_2
title: Bria 3.2
- local: api/pipelines/chroma
title: Chroma
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_flux
title: ControlNet with Flux.1
- local: api/pipelines/controlnet_hunyuandit
title: ControlNet with Hunyuan-DiT
- local: api/pipelines/controlnet_sd3
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/cosmos
title: Cosmos
- local: api/pipelines/ddim
title: DDIM
- local: api/pipelines/ddpm
title: DDPM
- local: api/pipelines/deepfloyd_if
title: DeepFloyd IF
- local: api/pipelines/diffedit
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/pix2pix
title: InstructPix2Pix
- local: api/pipelines/kandinsky
title: Kandinsky 2.1
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/kandinsky3
title: Kandinsky 3
- local: api/pipelines/kolors
title: Kolors
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
title: Paint by Example
- local: api/pipelines/pixart
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/qwenimage
title: QwenImage
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/stable_cascade
title: Stable Cascade
- sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-image
- local: api/pipelines/stable_diffusion/inpaint
title: Inpainting
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D
Upscaler
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
title: Stable Diffusion
- local: api/pipelines/stable_unclip
title: Stable unCLIP
- local: api/pipelines/unclip
title: unCLIP
- local: api/pipelines/unidiffuser
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/visualcloze
title: VisualCloze
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Image
- sections:
- local: api/pipelines/allegro
title: Allegro
- local: api/pipelines/cogvideox
title: CogVideoX
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/framepack
title: Framepack
- local: api/pipelines/hunyuan_video
title: HunyuanVideo
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/latte
title: Latte
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/mochi
title: Mochi
- local: api/pipelines/pia
title: Personalized Image Animator (PIA)
- local: api/pipelines/skyreels_v2
title: SkyReels-V2
- local: api/pipelines/stable_diffusion/svd
title: Image-to-video
- local: api/pipelines/stable_diffusion/inpaint
title: Inpainting
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
- local: api/pipelines/stable_unclip
title: Stable unCLIP
- local: api/pipelines/text_to_video
title: Text-to-video
- local: api/pipelines/text_to_video_zero
title: Text2Video-Zero
- local: api/pipelines/unclip
title: unCLIP
- local: api/pipelines/unidiffuser
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/visualcloze
title: VisualCloze
- local: api/pipelines/wan
title: Wan
- local: api/pipelines/wuerstchen
title: Wuerstchen
- title: Schedulers
sections:
title: Stable Video Diffusion
- local: api/pipelines/text_to_video
title: Text-to-video
- local: api/pipelines/text_to_video_zero
title: Text2Video-Zero
- local: api/pipelines/wan
title: Wan
title: Video
title: Pipelines
- sections:
- local: api/schedulers/overview
title: Overview
- local: api/schedulers/cm_stochastic_iterative
@@ -718,8 +711,8 @@
title: UniPCMultistepScheduler
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
- title: Internal classes
sections:
title: Schedulers
- sections:
- local: api/internal_classes_overview
title: Overview
- local: api/attnprocessor
@@ -736,3 +729,5 @@
title: VAE Image Processor
- local: api/video_processor
title: Video Processor
title: Internal classes
title: API
+3
View File
@@ -107,6 +107,9 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
[[autodoc]] loaders.lora_pipeline.QwenImageLoraLoaderMixin
## KandinskyLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.KandinskyLoraLoaderMixin
## LoraBaseMixin
[[autodoc]] loaders.lora_base.LoraBaseMixin
@@ -39,7 +39,7 @@ mask_url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images
original_image = load_image(img_url).resize((512, 512))
mask_image = load_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained("runwayml/stable-diffusion-inpainting")
pipe = StableDiffusionInpaintPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting")
pipe.vae = AsymmetricAutoencoderKL.from_pretrained("cross-attention/asymmetric-autoencoder-kl-x-1-5")
pipe.to("cuda")
@@ -21,7 +21,7 @@ The Stable Diffusion model can also infer depth based on an image using [MiDaS](
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionDepth2ImgPipeline
@@ -21,14 +21,14 @@ The Stable Diffusion model can also be applied to inpainting which lets you edit
## Tips
It is recommended to use this pipeline with checkpoints that have been specifically fine-tuned for inpainting, such
as [runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting). Default
as [stable-diffusion-v1-5/stable-diffusion-inpainting](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting). Default
text-to-image Stable Diffusion checkpoints, such as
[stable-diffusion-v1-5/stable-diffusion-v1-5](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) are also compatible but they might be less performant.
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionInpaintPipeline
@@ -17,7 +17,7 @@ The Stable Diffusion latent upscaler model was created by [Katherine Crowson](ht
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionLatentUpscalePipeline
@@ -22,7 +22,7 @@ Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B data
For more details about how Stable Diffusion works and how it differs from the base latent diffusion model, take a look at the Stability AI [announcement](https://stability.ai/blog/stable-diffusion-announcement) and our own [blog post](https://huggingface.co/blog/stable_diffusion#how-does-stable-diffusion-work) for more technical details.
You can find the original codebase for Stable Diffusion v1.0 at [CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion) and Stable Diffusion v2.0 at [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion) as well as their original scripts for various tasks. Additional official checkpoints for the different Stable Diffusion versions and tasks can be found on the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations. Explore these organizations to find the best checkpoint for your use-case!
You can find the original codebase for Stable Diffusion v1.0 at [CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion) and Stable Diffusion v2.0 at [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion) as well as their original scripts for various tasks. Additional official checkpoints for the different Stable Diffusion versions and tasks can be found on the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations. Explore these organizations to find the best checkpoint for your use-case!
The table below summarizes the available Stable Diffusion pipelines, their supported tasks, and an interactive demo:
@@ -64,7 +64,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<a href="./inpaint">StableDiffusionInpaint</a>
</td>
<td class="px-4 py-2 text-gray-700">inpainting</td>
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/runwayml/stable-diffusion-inpainting"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/stable-diffusion-v1-5/stable-diffusion-inpainting"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
@@ -36,7 +36,7 @@ Here are some examples for how to use Stable Diffusion 2 for each task:
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## Text-to-image
@@ -25,7 +25,7 @@ The abstract from the paper is:
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionPipeline
@@ -21,7 +21,7 @@ The Stable Diffusion upscaler diffusion model was created by the researchers and
> [!TIP]
> Make sure to check out the Stable Diffusion [Tips](overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
>
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis), [Runway](https://huggingface.co/runwayml), and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
> If you're interested in using one of the official checkpoints for a task, explore the [CompVis](https://huggingface.co/CompVis) and [Stability AI](https://huggingface.co/stabilityai) Hub organizations!
## StableDiffusionUpscalePipeline
+2 -2
View File
@@ -16,12 +16,12 @@ pipeline.unet.config["in_channels"]
4
```
Inpainting requires 9 channels in the input sample. You can check this value in a pretrained inpainting model like [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting):
Inpainting requires 9 channels in the input sample. You can check this value in a pretrained inpainting model like [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting):
```py
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", use_safetensors=True)
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting", use_safetensors=True)
pipeline.unet.config["in_channels"]
9
```
@@ -215,7 +215,7 @@ from diffusers import AutoPipelineForInpainting, LCMScheduler
from diffusers.utils import load_image, make_image_grid
pipe = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
variant="fp16",
).to("cuda")
+15 -15
View File
@@ -112,7 +112,7 @@ blurred_mask
## Popular models
[Stable Diffusion Inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2 Inpainting](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
[Stable Diffusion Inpainting](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting), [Stable Diffusion XL (SDXL) Inpainting](https://huggingface.co/diffusers/stable-diffusion-xl-1.0-inpainting-0.1), and [Kandinsky 2.2 Inpainting](https://huggingface.co/kandinsky-community/kandinsky-2-2-decoder-inpaint) are among the most popular models for inpainting. SDXL typically produces higher resolution images than Stable Diffusion v1.5, and Kandinsky 2.2 is also capable of generating high-quality images.
### Stable Diffusion Inpainting
@@ -124,7 +124,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -244,7 +244,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
```
</hfoption>
<hfoption id="runwayml/stable-diffusion-inpainting">
<hfoption id="stable-diffusion-v1-5/stable-diffusion-inpainting">
```py
import torch
@@ -252,7 +252,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -278,7 +278,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inpaint-specific.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">runwayml/stable-diffusion-inpainting</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">stable-diffusion-v1-5/stable-diffusion-inpainting</figcaption>
</div>
</div>
@@ -308,7 +308,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
```
</hfoption>
<hfoption id="runwayml/stable-diffusion-inpaint">
<hfoption id="stable-diffusion-v1-5/stable-diffusion-inpaint">
```py
import torch
@@ -316,7 +316,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -340,7 +340,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/specific-inpaint-basic.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">runwayml/stable-diffusion-inpainting</figcaption>
<figcaption class="mt-2 text-center text-sm text-gray-500">stable-diffusion-v1-5/stable-diffusion-inpainting</figcaption>
</div>
</div>
@@ -358,7 +358,7 @@ from diffusers.utils import load_image, make_image_grid
device = "cuda"
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
variant="fp16"
)
@@ -396,7 +396,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -441,7 +441,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -481,7 +481,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -606,7 +606,7 @@ from diffusers import AutoPipelineForInpainting, AutoPipelineForImage2Image
from diffusers.utils import load_image, make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -683,7 +683,7 @@ from diffusers import AutoPipelineForInpainting
from diffusers.utils import make_image_grid
pipeline = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16,
"stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16,
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
@@ -714,7 +714,7 @@ controlnet = ControlNetModel.from_pretrained("lllyasviel/control_v11p_sd15_inpai
# pass ControlNet to the pipeline
pipeline = StableDiffusionControlNetInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16"
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, torch_dtype=torch.float16, variant="fp16"
)
pipeline.enable_model_cpu_offload()
# remove following line if xFormers is not installed or you have PyTorch 2.0 or higher installed
+1 -1
View File
@@ -173,7 +173,7 @@ mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
path = "runwayml/stable-diffusion-inpainting"
path = "stable-diffusion-v1-5/stable-diffusion-inpainting"
run_compile = True # Set True / False
+2 -2
View File
@@ -28,12 +28,12 @@ pipeline.unet.config["in_channels"]
4
```
인페인팅은 입력 샘플에 9개의 채널이 필요합니다. [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting)와 같은 사전학습된 인페인팅 모델에서 이 값을 확인할 수 있습니다:
인페인팅은 입력 샘플에 9개의 채널이 필요합니다. [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting)와 같은 사전학습된 인페인팅 모델에서 이 값을 확인할 수 있습니다:
```py
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting")
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting")
pipeline.unet.config["in_channels"]
9
```
+2 -11
View File
@@ -14,7 +14,7 @@ specific language governing permissions and limitations under the License.
[[open-in-colab]]
[`StableDiffusionInpaintPipeline`]은 마스크와 텍스트 프롬프트를 제공하여 이미지의 특정 부분을 편집할 수 있도록 합니다. 이 기능은 인페인팅 작업을 위해 특별히 훈련된 [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting)과 같은 Stable Diffusion 버전을 사용합니다.
[`StableDiffusionInpaintPipeline`]은 마스크와 텍스트 프롬프트를 제공하여 이미지의 특정 부분을 편집할 수 있도록 합니다. 이 기능은 인페인팅 작업을 위해 특별히 훈련된 [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting)과 같은 Stable Diffusion 버전을 사용합니다.
먼저 [`StableDiffusionInpaintPipeline`] 인스턴스를 불러옵니다:
@@ -27,7 +27,7 @@ from io import BytesIO
from diffusers import StableDiffusionInpaintPipeline
pipeline = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
)
pipeline = pipeline.to("cuda")
@@ -61,12 +61,3 @@ image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
> [!WARNING]
> 이전의 실험적인 인페인팅 구현에서는 품질이 낮은 다른 프로세스를 사용했습니다. 이전 버전과의 호환성을 보장하기 위해 새 모델이 포함되지 않은 사전학습된 파이프라인을 불러오면 이전 인페인팅 방법이 계속 적용됩니다.
아래 Space에서 이미지 인페인팅을 직접 해보세요!
<iframe
src="https://runwayml-stable-diffusion-inpainting.hf.space"
frameborder="0"
width="850"
height="500"
></iframe>
+2 -2
View File
@@ -16,12 +16,12 @@ pipeline.unet.config["in_channels"]
4
```
而图像修复任务需要输入样本具有9个通道。您可以在 [`runwayml/stable-diffusion-inpainting`](https://huggingface.co/runwayml/stable-diffusion-inpainting) 这样的预训练修复模型中验证此参数:
而图像修复任务需要输入样本具有9个通道。您可以在 [`stable-diffusion-v1-5/stable-diffusion-inpainting`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-inpainting) 这样的预训练修复模型中验证此参数:
```python
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", use_safetensors=True)
pipeline = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-inpainting", use_safetensors=True)
pipeline.unet.config["in_channels"]
9
```
+1 -1
View File
@@ -1328,7 +1328,7 @@ model = CLIPSegForImageSegmentation.from_pretrained("CIDAS/clipseg-rd64-refined"
# Load Stable Diffusion Inpainting Pipeline with custom pipeline
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
custom_pipeline="text_inpainting",
segmentation_model=model,
segmentation_processor=processor
@@ -126,7 +126,7 @@ EXAMPLE_DOC_STRING = """
... "lllyasviel/control_v11p_sd15_inpaint", torch_dtype=torch.float16
... )
>>> pipe = StableDiffusionControlNetInpaintPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config)
@@ -347,7 +347,7 @@ class AdaptiveMaskInpaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -429,8 +429,8 @@ class AdaptiveMaskInpaintPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -970,7 +970,7 @@ class AdaptiveMaskInpaintPipeline(
>>> default_mask_image = download_image(mask_url).resize((512, 512))
>>> pipe = AdaptiveMaskInpaintPipeline.from_pretrained(
... "runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-inpainting", torch_dtype=torch.float16
... )
>>> pipe = pipe.to("cuda")
@@ -1095,7 +1095,7 @@ class AdaptiveMaskInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
@@ -62,7 +62,7 @@ class ComposableStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin)
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -145,8 +145,8 @@ class ComposableStableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin)
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -1276,7 +1276,7 @@ class FrescoV2VPipeline(StableDiffusionControlNetImg2ImgPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
+1 -1
View File
@@ -678,7 +678,7 @@ class StableDiffusionHDPainterPipeline(StableDiffusionInpaintPipeline):
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+1 -1
View File
@@ -78,7 +78,7 @@ class ImageToImageInpaintingPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+3 -3
View File
@@ -86,7 +86,7 @@ class InstaFlowPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -165,8 +165,8 @@ class InstaFlowPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+3 -3
View File
@@ -166,7 +166,7 @@ class IPAdapterFaceIDStableDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -247,8 +247,8 @@ class IPAdapterFaceIDStableDiffusionPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -414,7 +414,7 @@ class StableDiffusionHighResFixPipeline(StableDiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -222,7 +222,7 @@ class LatentConsistencyModelWalkPipeline(
supports [`LCMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
+3 -3
View File
@@ -302,7 +302,7 @@ class LLMGroundedDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -392,8 +392,8 @@ class LLMGroundedDiffusionPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+2 -2
View File
@@ -552,8 +552,8 @@ class StableDiffusionLongPromptWeightingPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1765,7 +1765,7 @@ class SDXLLongPromptWeightingPipeline(
# Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != num_channels_unet:
+2 -2
View File
@@ -3729,8 +3729,8 @@ class MatryoshkaPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -78,7 +78,7 @@ class MultilingualStableDiffusion(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -1607,7 +1607,7 @@ class KolorsControlNetInpaintPipeline(
# 9. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+3 -3
View File
@@ -135,7 +135,7 @@ class FabricPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
"""
@@ -163,8 +163,8 @@ class FabricPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1487,7 +1487,7 @@ class KolorsInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+3 -3
View File
@@ -106,7 +106,7 @@ class Prompt2PromptPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -187,8 +187,8 @@ class Prompt2PromptPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -1730,7 +1730,7 @@ class StyleAlignedSDXLPipeline(
# Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != num_channels_unet:
@@ -59,7 +59,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
@@ -392,7 +392,7 @@ class StableDiffusionBoxDiffPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -473,8 +473,8 @@ class StableDiffusionBoxDiffPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -42,7 +42,7 @@ EXAMPLE_DOC_STRING = """
```py
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
>>> image = pipe(prompt).images[0]
@@ -359,7 +359,7 @@ class StableDiffusionPAGPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -440,8 +440,8 @@ class StableDiffusionPAGPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -100,7 +100,7 @@ class StableDiffusionUpscaleLDM3DPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -2042,7 +2042,7 @@ class StableDiffusionXL_AE_Pipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
@@ -188,7 +188,7 @@ class StableDiffusionXLControlNetAdapterPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -330,7 +330,7 @@ class StableDiffusionXLControlNetAdapterInpaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
requires_aesthetics_score (`bool`, *optional*, defaults to `"False"`):
@@ -1569,7 +1569,7 @@ class StableDiffusionXLControlNetAdapterInpaintPipeline(
# 8. Check that sizes of mask, masked image and latents match
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
num_channels_mask = mask.shape[1]
num_channels_masked_image = masked_image_latents.shape[1]
if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels:
+4 -4
View File
@@ -46,7 +46,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = StableDiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> prompt = "a photo of an astronaut riding a horse on mars"
@@ -86,7 +86,7 @@ class Zero1to3StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
cc_projection ([`CCProjection`]):
@@ -164,8 +164,8 @@ class Zero1to3StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -288,7 +288,7 @@ class RerenderAVideoPipeline(StableDiffusionControlNetImg2ImgPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+1 -1
View File
@@ -54,7 +54,7 @@ EXAMPLE_DOC_STRING = """
>>> # load control net and stable diffusion v1-5
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> # speed up diffusion process with faster scheduler and memory optimization
@@ -158,7 +158,7 @@ EXAMPLE_DOC_STRING = """
>>> # load control net and stable diffusion v1-5
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> # speed up diffusion process with faster scheduler and memory optimization
@@ -64,7 +64,7 @@ class StableDiffusionPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
+1 -1
View File
@@ -114,7 +114,7 @@ class SdeDragPipeline(DiffusionPipeline):
>>> from diffusers import DDIMScheduler, DiffusionPipeline
>>> # Load the pipeline
>>> model_path = "runwayml/stable-diffusion-v1-5"
>>> model_path = "stable-diffusion-v1-5/stable-diffusion-v1-5"
>>> scheduler = DDIMScheduler.from_pretrained(model_path, subfolder="scheduler")
>>> pipe = DiffusionPipeline.from_pretrained(model_path, scheduler=scheduler, custom_pipeline="sde_drag")
>>> pipe.to('cuda')
@@ -46,7 +46,7 @@ class StableDiffusionComparisonPipeline(DiffusionPipeline, StableDiffusionMixin)
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionMegaSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -36,7 +36,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe_controlnet = StableDiffusionControlNetImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
controlnet=controlnet,
safety_checker=None,
torch_dtype=torch.float16
@@ -81,7 +81,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-seg", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
)
>>> pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
@@ -80,7 +80,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-seg", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetInpaintImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-inpainting", controlnet=controlnet, safety_checker=None, torch_dtype=torch.float16
)
>>> pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
@@ -37,7 +37,7 @@ EXAMPLE_DOC_STRING = """
>>> controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
>>> pipe = StableDiffusionControlNetReferencePipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
controlnet=controlnet,
safety_checker=None,
torch_dtype=torch.float16
+4 -4
View File
@@ -43,7 +43,7 @@ EXAMPLE_DOC_STRING = """
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_ipex")
>>> pipe = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", custom_pipeline="stable_diffusion_ipex")
>>> # For Float32
>>> pipe.prepare_for_ipex(prompt, dtype=torch.float32, height=512, width=512) #value of image height/width should be consistent with the pipeline inference
@@ -85,7 +85,7 @@ class StableDiffusionIPEXPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -161,8 +161,8 @@ class StableDiffusionIPEXPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -47,7 +47,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionMegaSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -46,7 +46,7 @@ EXAMPLE_DOC_STRING = """
>>> input_image = load_image("https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png")
>>> pipe = StableDiffusionReferencePipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
"stable-diffusion-v1-5/stable-diffusion-v1-5",
safety_checker=None,
torch_dtype=torch.float16
).to('cuda:0')
@@ -112,7 +112,7 @@ class StableDiffusionReferencePipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -194,8 +194,8 @@ class StableDiffusionReferencePipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -167,7 +167,7 @@ class StableDiffusionRepaintPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -249,8 +249,8 @@ class StableDiffusionRepaintPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -678,7 +678,7 @@ class TensorRTStableDiffusionImg2ImgPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -766,8 +766,8 @@ class TensorRTStableDiffusionImg2ImgPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -682,7 +682,7 @@ class TensorRTStableDiffusionInpaintPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -770,8 +770,8 @@ class TensorRTStableDiffusionInpaintPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -594,7 +594,7 @@ class TensorRTStableDiffusionPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -682,8 +682,8 @@ class TensorRTStableDiffusionPipeline(DiffusionPipeline):
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
+1 -1
View File
@@ -52,7 +52,7 @@ class TextInpainting(DiffusionPipeline, StableDiffusionMixin):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -1223,7 +1223,7 @@ class AnyTextPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -5,7 +5,7 @@ This script was added by @thedarkzeno .
Please note that this script is not actively maintained, you can open an issue and tag @thedarkzeno or @patil-suraj though.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export OUTPUT_DIR="path-to-save-model"
@@ -29,7 +29,7 @@ Prior-preservation is used to avoid overfitting and language-drift. Refer to the
According to the paper, it's recommended to generate `num_epochs * num_samples` images for prior-preservation. 200-300 works well for most cases.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -60,7 +60,7 @@ With the help of gradient checkpointing and the 8-bit optimizer from bitsandbyte
To install `bitandbytes` please refer to this [readme](https://github.com/TimDettmers/bitsandbytes#requirements--installation).
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -92,7 +92,7 @@ Pass the `--train_text_encoder` argument to the script to enable training `text_
___Note: Training text encoder requires more memory, with this option the training won't fit on 16GB GPU. It needs at least 24GB VRAM.___
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
@@ -55,7 +55,7 @@ The Accelerate launch command is used to train a model using multiple GPUs and m
```
accelerate launch --mixed_precision "fp16" \
tutorial_train_ip-adapter.py \
--pretrained_model_name_or_path="runwayml/stable-diffusion-v1-5/" \
--pretrained_model_name_or_path="stable-diffusion-v1-5/stable-diffusion-v1-5/" \
--image_encoder_path="{image_encoder_path}" \
--data_json_file="{data.json}" \
--data_root_path="{image_path}" \
@@ -73,7 +73,7 @@ tutorial_train_ip-adapter.py \
```
accelerate launch --num_processes 8 --multi_gpu --mixed_precision "fp16" \
tutorial_train_ip-adapter.py \
--pretrained_model_name_or_path="runwayml/stable-diffusion-v1-5/" \
--pretrained_model_name_or_path="stable-diffusion-v1-5/stable-diffusion-v1-5/" \
--image_encoder_path="{image_encoder_path}" \
--data_json_file="{data.json}" \
--data_root_path="{image_path}" \
@@ -27,7 +27,7 @@ You can build multiple datasets for every subject and upload them to the 🤗 hu
Before launching the training script, make sure to select the inpainting the target model, the output directory and the 🤗 datasets.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export MODEL_NAME="stable-diffusion-v1-5/stable-diffusion-inpainting"
export OUTPUT_DIR="path-to-save-model"
export DATASET_1="gzguevara/mr_potato_head_masked"
@@ -177,7 +177,7 @@ class PromptDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for more details
about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
@@ -238,7 +238,7 @@ def parse_args() -> argparse.Namespace:
# EXAMPLE USAGE:
#
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "runwayml/stable-diffusion-v1-5" --subfolder "vae" --input_image "foo.png"
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "stable-diffusion-v1-5/stable-diffusion-v1-5" --subfolder "vae" --input_image "foo.png"
#
# python vae_roundtrip.py --use_cuda --pretrained_model_name_or_path "madebyollin/taesd" --use_tiny_nn --input_image "foo.png"
#
+7 -4
View File
@@ -24,7 +24,8 @@ args = args.parse_args()
def _extract_into_tensor(arr, timesteps, broadcast_shape):
# from: https://github.com/openai/guided-diffusion/blob/22e0df8183507e13a7813f8d38d51b072ca1e67c/guided_diffusion/gaussian_diffusion.py#L895 """
# from: https://github.com/openai/guided-diffusion/blob/22e0df8183507e13a7813f8d38d51b072ca1e67c/guided_diffusion/gaussian_diffusion.py#L895
# """
res = arr[timesteps].float()
dims_to_append = len(broadcast_shape) - len(res.shape)
return res[(...,) + (None,) * dims_to_append]
@@ -507,7 +508,9 @@ def rename_state_dict(sd, embedding):
# encode with stable diffusion vae
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16
)
pipe.vae.cuda()
# construct original decoder with jitted model
@@ -1090,7 +1093,7 @@ def new_constructor(self, **kwargs):
Encoder.__init__ = new_constructor
vae = AutoencoderKL.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="vae")
vae = AutoencoderKL.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="vae")
consistency_vae = ConsistencyDecoderVAE(
encoder_args=vae.encoder.constructor_arguments,
decoder_args=unet.config,
@@ -1117,7 +1120,7 @@ print((sample_consistency_orig - sample_consistency_new_3).abs().sum())
print("running with diffusers pipeline")
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", vae=consistency_vae, torch_dtype=torch.float16
"stable-diffusion-v1-5/stable-diffusion-v1-5", vae=consistency_vae, torch_dtype=torch.float16
)
pipe.to("cuda")
+4
View File
@@ -220,6 +220,7 @@ else:
"HunyuanVideoTransformer3DModel",
"I2VGenXLUNet",
"Kandinsky3UNet",
"Kandinsky5Transformer3DModel",
"LatteTransformer3DModel",
"LTXVideoTransformer3DModel",
"Lumina2Transformer2DModel",
@@ -474,6 +475,7 @@ else:
"ImageTextPipelineOutput",
"Kandinsky3Img2ImgPipeline",
"Kandinsky3Pipeline",
"Kandinsky5T2VPipeline",
"KandinskyCombinedPipeline",
"KandinskyImg2ImgCombinedPipeline",
"KandinskyImg2ImgPipeline",
@@ -912,6 +914,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
HunyuanVideoTransformer3DModel,
I2VGenXLUNet,
Kandinsky3UNet,
Kandinsky5Transformer3DModel,
LatteTransformer3DModel,
LTXVideoTransformer3DModel,
Lumina2Transformer2DModel,
@@ -1136,6 +1139,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
ImageTextPipelineOutput,
Kandinsky3Img2ImgPipeline,
Kandinsky3Pipeline,
Kandinsky5T2VPipeline,
KandinskyCombinedPipeline,
KandinskyImg2ImgCombinedPipeline,
KandinskyImg2ImgPipeline,
+2
View File
@@ -77,6 +77,7 @@ if is_torch_available():
"SanaLoraLoaderMixin",
"Lumina2LoraLoaderMixin",
"WanLoraLoaderMixin",
"KandinskyLoraLoaderMixin",
"HiDreamImageLoraLoaderMixin",
"SkyReelsV2LoraLoaderMixin",
"QwenImageLoraLoaderMixin",
@@ -115,6 +116,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
FluxLoraLoaderMixin,
HiDreamImageLoraLoaderMixin,
HunyuanVideoLoraLoaderMixin,
KandinskyLoraLoaderMixin,
LoraLoaderMixin,
LTXVideoLoraLoaderMixin,
Lumina2LoraLoaderMixin,
+285
View File
@@ -3639,6 +3639,291 @@ class Lumina2LoraLoaderMixin(LoraBaseMixin):
super().unfuse_lora(components=components, **kwargs)
class KandinskyLoraLoaderMixin(LoraBaseMixin):
r"""
Load LoRA layers into [`Kandinsky5Transformer3DModel`],
"""
_lora_loadable_modules = ["transformer"]
transformer_name = TRANSFORMER_NAME
@classmethod
@validate_hf_hub_args
def lora_state_dict(
cls,
pretrained_model_name_or_path_or_dict: Union[str, Dict[str, torch.Tensor]],
**kwargs,
):
r"""
Return state dict for lora weights and the network alphas.
Parameters:
pretrained_model_name_or_path_or_dict (`str` or `os.PathLike` or `dict`):
Can be either:
- A string, the *model id* of a pretrained model hosted on the Hub.
- A path to a *directory* containing the model weights.
- A [torch state
dict](https://pytorch.org/tutorials/beginner/saving_loading_models.html#what-is-a-state-dict).
cache_dir (`Union[str, os.PathLike]`, *optional*):
Path to a directory where a downloaded pretrained model configuration is cached.
force_download (`bool`, *optional*, defaults to `False`):
Whether or not to force the (re-)download of the model weights.
proxies (`Dict[str, str]`, *optional*):
A dictionary of proxy servers to use by protocol or endpoint.
local_files_only (`bool`, *optional*, defaults to `False`):
Whether to only load local model weights and configuration files.
token (`str` or *bool*, *optional*):
The token to use as HTTP bearer authorization for remote files.
revision (`str`, *optional*, defaults to `"main"`):
The specific model version to use.
subfolder (`str`, *optional*, defaults to `""`):
The subfolder location of a model file within a larger model repository.
weight_name (`str`, *optional*, defaults to None):
Name of the serialized state dict file.
use_safetensors (`bool`, *optional*):
Whether to use safetensors for loading.
return_lora_metadata (`bool`, *optional*, defaults to False):
When enabled, additionally return the LoRA adapter metadata.
"""
# Load the main state dict first which has the LoRA layers
cache_dir = kwargs.pop("cache_dir", None)
force_download = kwargs.pop("force_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", None)
token = kwargs.pop("token", None)
revision = kwargs.pop("revision", None)
subfolder = kwargs.pop("subfolder", None)
weight_name = kwargs.pop("weight_name", None)
use_safetensors = kwargs.pop("use_safetensors", None)
return_lora_metadata = kwargs.pop("return_lora_metadata", False)
allow_pickle = False
if use_safetensors is None:
use_safetensors = True
allow_pickle = True
user_agent = {"file_type": "attn_procs_weights", "framework": "pytorch"}
state_dict, metadata = _fetch_state_dict(
pretrained_model_name_or_path_or_dict=pretrained_model_name_or_path_or_dict,
weight_name=weight_name,
use_safetensors=use_safetensors,
local_files_only=local_files_only,
cache_dir=cache_dir,
force_download=force_download,
proxies=proxies,
token=token,
revision=revision,
subfolder=subfolder,
user_agent=user_agent,
allow_pickle=allow_pickle,
)
is_dora_scale_present = any("dora_scale" in k for k in state_dict)
if is_dora_scale_present:
warn_msg = "It seems like you are using a DoRA checkpoint that is not compatible in Diffusers at the moment. So, we are going to filter out the keys associated to 'dora_scale` from the state dict. If you think this is a mistake please open an issue https://github.com/huggingface/diffusers/issues/new."
logger.warning(warn_msg)
state_dict = {k: v for k, v in state_dict.items() if "dora_scale" not in k}
out = (state_dict, metadata) if return_lora_metadata else state_dict
return out
def load_lora_weights(
self,
pretrained_model_name_or_path_or_dict: Union[str, Dict[str, torch.Tensor]],
adapter_name: Optional[str] = None,
hotswap: bool = False,
**kwargs,
):
"""
Load LoRA weights specified in `pretrained_model_name_or_path_or_dict` into `self.transformer`
Parameters:
pretrained_model_name_or_path_or_dict (`str` or `os.PathLike` or `dict`):
See [`~loaders.KandinskyLoraLoaderMixin.lora_state_dict`].
adapter_name (`str`, *optional*):
Adapter name to be used for referencing the loaded adapter model.
hotswap (`bool`, *optional*):
Whether to substitute an existing (LoRA) adapter with the newly loaded adapter in-place.
low_cpu_mem_usage (`bool`, *optional*):
Speed up model loading by only loading the pretrained LoRA weights and not initializing the random
weights.
kwargs (`dict`, *optional*):
See [`~loaders.KandinskyLoraLoaderMixin.lora_state_dict`].
"""
if not USE_PEFT_BACKEND:
raise ValueError("PEFT backend is required for this method.")
low_cpu_mem_usage = kwargs.pop("low_cpu_mem_usage", _LOW_CPU_MEM_USAGE_DEFAULT_LORA)
if low_cpu_mem_usage and not is_peft_version(">=", "0.13.1"):
raise ValueError(
"`low_cpu_mem_usage=True` is not compatible with this `peft` version. Please update it with `pip install -U peft`."
)
# if a dict is passed, copy it instead of modifying it inplace
if isinstance(pretrained_model_name_or_path_or_dict, dict):
pretrained_model_name_or_path_or_dict = pretrained_model_name_or_path_or_dict.copy()
# First, ensure that the checkpoint is a compatible one and can be successfully loaded.
kwargs["return_lora_metadata"] = True
state_dict, metadata = self.lora_state_dict(pretrained_model_name_or_path_or_dict, **kwargs)
is_correct_format = all("lora" in key for key in state_dict.keys())
if not is_correct_format:
raise ValueError("Invalid LoRA checkpoint.")
# Load LoRA into transformer
self.load_lora_into_transformer(
state_dict,
transformer=getattr(self, self.transformer_name) if not hasattr(self, "transformer") else self.transformer,
adapter_name=adapter_name,
metadata=metadata,
_pipeline=self,
low_cpu_mem_usage=low_cpu_mem_usage,
hotswap=hotswap,
)
@classmethod
def load_lora_into_transformer(
cls,
state_dict,
transformer,
adapter_name=None,
_pipeline=None,
low_cpu_mem_usage=False,
hotswap: bool = False,
metadata=None,
):
"""
Load the LoRA layers specified in `state_dict` into `transformer`.
Parameters:
state_dict (`dict`):
A standard state dict containing the lora layer parameters.
transformer (`Kandinsky5Transformer3DModel`):
The transformer model to load the LoRA layers into.
adapter_name (`str`, *optional*):
Adapter name to be used for referencing the loaded adapter model.
low_cpu_mem_usage (`bool`, *optional*):
Speed up model loading by only loading the pretrained LoRA weights.
hotswap (`bool`, *optional*):
See [`~loaders.KandinskyLoraLoaderMixin.load_lora_weights`].
metadata (`dict`):
Optional LoRA adapter metadata.
"""
if low_cpu_mem_usage and not is_peft_version(">=", "0.13.1"):
raise ValueError(
"`low_cpu_mem_usage=True` is not compatible with this `peft` version. Please update it with `pip install -U peft`."
)
# Load the layers corresponding to transformer.
logger.info(f"Loading {cls.transformer_name}.")
transformer.load_lora_adapter(
state_dict,
network_alphas=None,
adapter_name=adapter_name,
metadata=metadata,
_pipeline=_pipeline,
low_cpu_mem_usage=low_cpu_mem_usage,
hotswap=hotswap,
)
@classmethod
def save_lora_weights(
cls,
save_directory: Union[str, os.PathLike],
transformer_lora_layers: Dict[str, Union[torch.nn.Module, torch.Tensor]] = None,
is_main_process: bool = True,
weight_name: str = None,
save_function: Callable = None,
safe_serialization: bool = True,
transformer_lora_adapter_metadata=None,
):
r"""
Save the LoRA parameters corresponding to the transformer and text encoders.
Arguments:
save_directory (`str` or `os.PathLike`):
Directory to save LoRA parameters to.
transformer_lora_layers (`Dict[str, torch.nn.Module]` or `Dict[str, torch.Tensor]`):
State dict of the LoRA layers corresponding to the `transformer`.
is_main_process (`bool`, *optional*, defaults to `True`):
Whether the process calling this is the main process.
save_function (`Callable`):
The function to use to save the state dictionary.
safe_serialization (`bool`, *optional*, defaults to `True`):
Whether to save the model using `safetensors` or the traditional PyTorch way.
transformer_lora_adapter_metadata:
LoRA adapter metadata associated with the transformer.
"""
lora_layers = {}
lora_metadata = {}
if transformer_lora_layers:
lora_layers[cls.transformer_name] = transformer_lora_layers
lora_metadata[cls.transformer_name] = transformer_lora_adapter_metadata
if not lora_layers:
raise ValueError("You must pass at least one of `transformer_lora_layers`")
cls._save_lora_weights(
save_directory=save_directory,
lora_layers=lora_layers,
lora_metadata=lora_metadata,
is_main_process=is_main_process,
weight_name=weight_name,
save_function=save_function,
safe_serialization=safe_serialization,
)
def fuse_lora(
self,
components: List[str] = ["transformer"],
lora_scale: float = 1.0,
safe_fusing: bool = False,
adapter_names: Optional[List[str]] = None,
**kwargs,
):
r"""
Fuses the LoRA parameters into the original parameters of the corresponding blocks.
Args:
components: (`List[str]`): List of LoRA-injectable components to fuse the LoRAs into.
lora_scale (`float`, defaults to 1.0):
Controls how much to influence the outputs with the LoRA parameters.
safe_fusing (`bool`, defaults to `False`):
Whether to check fused weights for NaN values before fusing.
adapter_names (`List[str]`, *optional*):
Adapter names to be used for fusing.
Example:
```py
from diffusers import Kandinsky5T2VPipeline
pipeline = Kandinsky5T2VPipeline.from_pretrained("ai-forever/Kandinsky-5.0-T2V")
pipeline.load_lora_weights("path/to/lora.safetensors")
pipeline.fuse_lora(lora_scale=0.7)
```
"""
super().fuse_lora(
components=components,
lora_scale=lora_scale,
safe_fusing=safe_fusing,
adapter_names=adapter_names,
**kwargs,
)
def unfuse_lora(self, components: List[str] = ["transformer"], **kwargs):
r"""
Reverses the effect of [`pipe.fuse_lora()`].
Args:
components (`List[str]`): List of LoRA-injectable components to unfuse LoRA from.
"""
super().unfuse_lora(components=components, **kwargs)
class WanLoraLoaderMixin(LoraBaseMixin):
r"""
Load LoRA layers into [`WanTransformer3DModel`]. Specific to [`WanPipeline`] and `[WanImageToVideoPipeline`].
+1 -1
View File
@@ -293,7 +293,7 @@ class PeftAdapterMixin:
# For hotswapping, we need the adapter name to be present in the state dict keys
new_sd = {}
for k, v in sd.items():
if k.endswith("lora_A.weight") or key.endswith("lora_B.weight"):
if k.endswith("lora_A.weight") or k.endswith("lora_B.weight"):
k = k[: -len(".weight")] + f".{adapter_name}.weight"
elif k.endswith("lora_B.bias"): # lora_bias=True option
k = k[: -len(".bias")] + f".{adapter_name}.bias"
+2
View File
@@ -91,6 +91,7 @@ if is_torch_available():
_import_structure["transformers.transformer_hidream_image"] = ["HiDreamImageTransformer2DModel"]
_import_structure["transformers.transformer_hunyuan_video"] = ["HunyuanVideoTransformer3DModel"]
_import_structure["transformers.transformer_hunyuan_video_framepack"] = ["HunyuanVideoFramepackTransformer3DModel"]
_import_structure["transformers.transformer_kandinsky"] = ["Kandinsky5Transformer3DModel"]
_import_structure["transformers.transformer_ltx"] = ["LTXVideoTransformer3DModel"]
_import_structure["transformers.transformer_lumina2"] = ["Lumina2Transformer2DModel"]
_import_structure["transformers.transformer_mochi"] = ["MochiTransformer3DModel"]
@@ -182,6 +183,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
HunyuanDiT2DModel,
HunyuanVideoFramepackTransformer3DModel,
HunyuanVideoTransformer3DModel,
Kandinsky5Transformer3DModel,
LatteTransformer3DModel,
LTXVideoTransformer3DModel,
Lumina2Transformer2DModel,
+2 -2
View File
@@ -128,13 +128,13 @@ class AutoModel(ConfigMixin):
```py
from diffusers import AutoModel
unet = AutoModel.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="unet")
unet = AutoModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="unet")
```
If you get the error message below, you need to finetune the weights for your downstream task:
```bash
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at runwayml/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at stable-diffusion-v1-5/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
- conv_in.weight: found shape torch.Size([320, 4, 3, 3]) in the checkpoint and torch.Size([320, 9, 3, 3]) in the model instantiated
You should probably TRAIN this model on a down-stream task to be able to use it for predictions and inference.
```
+9 -9
View File
@@ -113,14 +113,14 @@ class FlaxModelMixin(PushToHubMixin):
>>> from diffusers import FlaxUNet2DConditionModel
>>> # load model
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> # By default, the model parameters will be in fp32 precision, to cast these to bfloat16 precision
>>> params = model.to_bf16(params)
>>> # If you don't want to cast certain parameters (for example layer norm bias and scale)
>>> # then pass the mask as follows
>>> from flax import traverse_util
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> flat_params = traverse_util.flatten_dict(params)
>>> mask = {
... path: (path[-2] != ("LayerNorm", "bias") and path[-2:] != ("LayerNorm", "scale"))
@@ -149,7 +149,7 @@ class FlaxModelMixin(PushToHubMixin):
>>> from diffusers import FlaxUNet2DConditionModel
>>> # Download model and configuration from huggingface.co
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> # By default, the model params will be in fp32, to illustrate the use of this method,
>>> # we'll first cast to fp16 and back to fp32
>>> params = model.to_f16(params)
@@ -179,14 +179,14 @@ class FlaxModelMixin(PushToHubMixin):
>>> from diffusers import FlaxUNet2DConditionModel
>>> # load model
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> # By default, the model params will be in fp32, to cast these to float16
>>> params = model.to_fp16(params)
>>> # If you want don't want to cast certain parameters (for example layer norm bias and scale)
>>> # then pass the mask as follows
>>> from flax import traverse_util
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> flat_params = traverse_util.flatten_dict(params)
>>> mask = {
... path: (path[-2] != ("LayerNorm", "bias") and path[-2:] != ("LayerNorm", "scale"))
@@ -216,8 +216,8 @@ class FlaxModelMixin(PushToHubMixin):
pretrained_model_name_or_path (`str` or `os.PathLike`):
Can be either:
- A string, the *model id* (for example `runwayml/stable-diffusion-v1-5`) of a pretrained model
hosted on the Hub.
- A string, the *model id* (for example `stable-diffusion-v1-5/stable-diffusion-v1-5`) of a
pretrained model hosted on the Hub.
- A path to a *directory* (for example `./my_model_directory`) containing the model weights saved
using [`~FlaxModelMixin.save_pretrained`].
dtype (`jax.numpy.dtype`, *optional*, defaults to `jax.numpy.float32`):
@@ -271,7 +271,7 @@ class FlaxModelMixin(PushToHubMixin):
>>> from diffusers import FlaxUNet2DConditionModel
>>> # Download model and configuration from huggingface.co and cache.
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
>>> # Model was saved using *save_pretrained('./test/saved_model/')* (for example purposes, not runnable).
>>> model, params = FlaxUNet2DConditionModel.from_pretrained("./test/saved_model/")
```
@@ -279,7 +279,7 @@ class FlaxModelMixin(PushToHubMixin):
If you get the error message below, you need to finetune the weights for your downstream task:
```bash
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at runwayml/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at stable-diffusion-v1-5/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
- conv_in.weight: found shape torch.Size([320, 4, 3, 3]) in the checkpoint and torch.Size([320, 9, 3, 3]) in the model instantiated
You should probably TRAIN this model on a down-stream task to be able to use it for predictions and inference.
```
+3 -3
View File
@@ -923,13 +923,13 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
```py
from diffusers import UNet2DConditionModel
unet = UNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="unet")
unet = UNet2DConditionModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="unet")
```
If you get the error message below, you need to finetune the weights for your downstream task:
```bash
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at runwayml/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
Some weights of UNet2DConditionModel were not initialized from the model checkpoint at stable-diffusion-v1-5/stable-diffusion-v1-5 and are newly initialized because the shapes did not match:
- conv_in.weight: found shape torch.Size([320, 4, 3, 3]) in the checkpoint and torch.Size([320, 9, 3, 3]) in the model instantiated
You should probably TRAIN this model on a down-stream task to be able to use it for predictions and inference.
```
@@ -1800,7 +1800,7 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
```py
from diffusers import UNet2DConditionModel
model_id = "runwayml/stable-diffusion-v1-5"
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
unet = UNet2DConditionModel.from_pretrained(model_id, subfolder="unet")
unet.num_parameters(only_trainable=True)
859520964
@@ -27,6 +27,7 @@ if is_torch_available():
from .transformer_hidream_image import HiDreamImageTransformer2DModel
from .transformer_hunyuan_video import HunyuanVideoTransformer3DModel
from .transformer_hunyuan_video_framepack import HunyuanVideoFramepackTransformer3DModel
from .transformer_kandinsky import Kandinsky5Transformer3DModel
from .transformer_ltx import LTXVideoTransformer3DModel
from .transformer_lumina2 import Lumina2Transformer2DModel
from .transformer_mochi import MochiTransformer3DModel
@@ -0,0 +1,667 @@
# Copyright 2025 The Kandinsky Team and The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import inspect
import math
from typing import Any, Dict, Optional, Tuple, Union
import torch
import torch.nn as nn
import torch.nn.functional as F
from torch import Tensor
from ...configuration_utils import ConfigMixin, register_to_config
from ...loaders import FromOriginalModelMixin, PeftAdapterMixin
from ...utils import (
logging,
)
from ..attention import AttentionMixin, AttentionModuleMixin
from ..attention_dispatch import _CAN_USE_FLEX_ATTN, dispatch_attention_fn
from ..cache_utils import CacheMixin
from ..modeling_outputs import Transformer2DModelOutput
from ..modeling_utils import ModelMixin
logger = logging.get_logger(__name__)
def get_freqs(dim, max_period=10000.0):
freqs = torch.exp(-math.log(max_period) * torch.arange(start=0, end=dim, dtype=torch.float32) / dim)
return freqs
def fractal_flatten(x, rope, shape, block_mask=False):
if block_mask:
pixel_size = 8
x = local_patching(x, shape, (1, pixel_size, pixel_size), dim=1)
rope = local_patching(rope, shape, (1, pixel_size, pixel_size), dim=1)
x = x.flatten(1, 2)
rope = rope.flatten(1, 2)
else:
x = x.flatten(1, 3)
rope = rope.flatten(1, 3)
return x, rope
def fractal_unflatten(x, shape, block_mask=False):
if block_mask:
pixel_size = 8
x = x.reshape(x.shape[0], -1, pixel_size**2, *x.shape[2:])
x = local_merge(x, shape, (1, pixel_size, pixel_size), dim=1)
else:
x = x.reshape(*shape, *x.shape[2:])
return x
def local_patching(x, shape, group_size, dim=0):
batch_size, duration, height, width = shape
g1, g2, g3 = group_size
x = x.reshape(
*x.shape[:dim],
duration // g1,
g1,
height // g2,
g2,
width // g3,
g3,
*x.shape[dim + 3 :],
)
x = x.permute(
*range(len(x.shape[:dim])),
dim,
dim + 2,
dim + 4,
dim + 1,
dim + 3,
dim + 5,
*range(dim + 6, len(x.shape)),
)
x = x.flatten(dim, dim + 2).flatten(dim + 1, dim + 3)
return x
def local_merge(x, shape, group_size, dim=0):
batch_size, duration, height, width = shape
g1, g2, g3 = group_size
x = x.reshape(
*x.shape[:dim],
duration // g1,
height // g2,
width // g3,
g1,
g2,
g3,
*x.shape[dim + 2 :],
)
x = x.permute(
*range(len(x.shape[:dim])),
dim,
dim + 3,
dim + 1,
dim + 4,
dim + 2,
dim + 5,
*range(dim + 6, len(x.shape)),
)
x = x.flatten(dim, dim + 1).flatten(dim + 1, dim + 2).flatten(dim + 2, dim + 3)
return x
def nablaT_v2(
q: Tensor,
k: Tensor,
sta: Tensor,
thr: float = 0.9,
):
if _CAN_USE_FLEX_ATTN:
from torch.nn.attention.flex_attention import BlockMask
else:
raise ValueError("Nabla attention is not supported with this version of PyTorch")
q = q.transpose(1, 2).contiguous()
k = k.transpose(1, 2).contiguous()
# Map estimation
B, h, S, D = q.shape
s1 = S // 64
qa = q.reshape(B, h, s1, 64, D).mean(-2)
ka = k.reshape(B, h, s1, 64, D).mean(-2).transpose(-2, -1)
map = qa @ ka
map = torch.softmax(map / math.sqrt(D), dim=-1)
# Map binarization
vals, inds = map.sort(-1)
cvals = vals.cumsum_(-1)
mask = (cvals >= 1 - thr).int()
mask = mask.gather(-1, inds.argsort(-1))
mask = torch.logical_or(mask, sta)
# BlockMask creation
kv_nb = mask.sum(-1).to(torch.int32)
kv_inds = mask.argsort(dim=-1, descending=True).to(torch.int32)
return BlockMask.from_kv_blocks(torch.zeros_like(kv_nb), kv_inds, kv_nb, kv_inds, BLOCK_SIZE=64, mask_mod=None)
class Kandinsky5TimeEmbeddings(nn.Module):
def __init__(self, model_dim, time_dim, max_period=10000.0):
super().__init__()
assert model_dim % 2 == 0
self.model_dim = model_dim
self.max_period = max_period
self.freqs = get_freqs(self.model_dim // 2, self.max_period)
self.in_layer = nn.Linear(model_dim, time_dim, bias=True)
self.activation = nn.SiLU()
self.out_layer = nn.Linear(time_dim, time_dim, bias=True)
@torch.autocast(device_type="cuda", dtype=torch.float32)
def forward(self, time):
args = torch.outer(time, self.freqs.to(device=time.device))
time_embed = torch.cat([torch.cos(args), torch.sin(args)], dim=-1)
time_embed = self.out_layer(self.activation(self.in_layer(time_embed)))
return time_embed
class Kandinsky5TextEmbeddings(nn.Module):
def __init__(self, text_dim, model_dim):
super().__init__()
self.in_layer = nn.Linear(text_dim, model_dim, bias=True)
self.norm = nn.LayerNorm(model_dim, elementwise_affine=True)
def forward(self, text_embed):
text_embed = self.in_layer(text_embed)
return self.norm(text_embed).type_as(text_embed)
class Kandinsky5VisualEmbeddings(nn.Module):
def __init__(self, visual_dim, model_dim, patch_size):
super().__init__()
self.patch_size = patch_size
self.in_layer = nn.Linear(math.prod(patch_size) * visual_dim, model_dim)
def forward(self, x):
batch_size, duration, height, width, dim = x.shape
x = (
x.view(
batch_size,
duration // self.patch_size[0],
self.patch_size[0],
height // self.patch_size[1],
self.patch_size[1],
width // self.patch_size[2],
self.patch_size[2],
dim,
)
.permute(0, 1, 3, 5, 2, 4, 6, 7)
.flatten(4, 7)
)
return self.in_layer(x)
class Kandinsky5RoPE1D(nn.Module):
def __init__(self, dim, max_pos=1024, max_period=10000.0):
super().__init__()
self.max_period = max_period
self.dim = dim
self.max_pos = max_pos
freq = get_freqs(dim // 2, max_period)
pos = torch.arange(max_pos, dtype=freq.dtype)
self.register_buffer("args", torch.outer(pos, freq), persistent=False)
def forward(self, pos):
args = self.args[pos]
cosine = torch.cos(args)
sine = torch.sin(args)
rope = torch.stack([cosine, -sine, sine, cosine], dim=-1)
rope = rope.view(*rope.shape[:-1], 2, 2)
return rope.unsqueeze(-4)
class Kandinsky5RoPE3D(nn.Module):
def __init__(self, axes_dims, max_pos=(128, 128, 128), max_period=10000.0):
super().__init__()
self.axes_dims = axes_dims
self.max_pos = max_pos
self.max_period = max_period
for i, (axes_dim, ax_max_pos) in enumerate(zip(axes_dims, max_pos)):
freq = get_freqs(axes_dim // 2, max_period)
pos = torch.arange(ax_max_pos, dtype=freq.dtype)
self.register_buffer(f"args_{i}", torch.outer(pos, freq), persistent=False)
def forward(self, shape, pos, scale_factor=(1.0, 1.0, 1.0)):
batch_size, duration, height, width = shape
args_t = self.args_0[pos[0]] / scale_factor[0]
args_h = self.args_1[pos[1]] / scale_factor[1]
args_w = self.args_2[pos[2]] / scale_factor[2]
args = torch.cat(
[
args_t.view(1, duration, 1, 1, -1).repeat(batch_size, 1, height, width, 1),
args_h.view(1, 1, height, 1, -1).repeat(batch_size, duration, 1, width, 1),
args_w.view(1, 1, 1, width, -1).repeat(batch_size, duration, height, 1, 1),
],
dim=-1,
)
cosine = torch.cos(args)
sine = torch.sin(args)
rope = torch.stack([cosine, -sine, sine, cosine], dim=-1)
rope = rope.view(*rope.shape[:-1], 2, 2)
return rope.unsqueeze(-4)
class Kandinsky5Modulation(nn.Module):
def __init__(self, time_dim, model_dim, num_params):
super().__init__()
self.activation = nn.SiLU()
self.out_layer = nn.Linear(time_dim, num_params * model_dim)
self.out_layer.weight.data.zero_()
self.out_layer.bias.data.zero_()
@torch.autocast(device_type="cuda", dtype=torch.float32)
def forward(self, x):
return self.out_layer(self.activation(x))
class Kandinsky5AttnProcessor:
_attention_backend = None
_parallel_config = None
def __init__(self):
if not hasattr(F, "scaled_dot_product_attention"):
raise ImportError(f"{self.__class__.__name__} requires PyTorch 2.0. Please upgrade your pytorch version.")
def __call__(self, attn, hidden_states, encoder_hidden_states=None, rotary_emb=None, sparse_params=None):
# query, key, value = self.get_qkv(x)
query = attn.to_query(hidden_states)
if encoder_hidden_states is not None:
key = attn.to_key(encoder_hidden_states)
value = attn.to_value(encoder_hidden_states)
shape, cond_shape = query.shape[:-1], key.shape[:-1]
query = query.reshape(*shape, attn.num_heads, -1)
key = key.reshape(*cond_shape, attn.num_heads, -1)
value = value.reshape(*cond_shape, attn.num_heads, -1)
else:
key = attn.to_key(hidden_states)
value = attn.to_value(hidden_states)
shape = query.shape[:-1]
query = query.reshape(*shape, attn.num_heads, -1)
key = key.reshape(*shape, attn.num_heads, -1)
value = value.reshape(*shape, attn.num_heads, -1)
# query, key = self.norm_qk(query, key)
query = attn.query_norm(query.float()).type_as(query)
key = attn.key_norm(key.float()).type_as(key)
def apply_rotary(x, rope):
x_ = x.reshape(*x.shape[:-1], -1, 1, 2).to(torch.float32)
x_out = (rope * x_).sum(dim=-1)
return x_out.reshape(*x.shape).to(torch.bfloat16)
if rotary_emb is not None:
query = apply_rotary(query, rotary_emb).type_as(query)
key = apply_rotary(key, rotary_emb).type_as(key)
if sparse_params is not None:
attn_mask = nablaT_v2(
query,
key,
sparse_params["sta_mask"],
thr=sparse_params["P"],
)
else:
attn_mask = None
hidden_states = dispatch_attention_fn(
query,
key,
value,
attn_mask=attn_mask,
backend=self._attention_backend,
parallel_config=self._parallel_config,
)
hidden_states = hidden_states.flatten(-2, -1)
attn_out = attn.out_layer(hidden_states)
return attn_out
class Kandinsky5Attention(nn.Module, AttentionModuleMixin):
_default_processor_cls = Kandinsky5AttnProcessor
_available_processors = [
Kandinsky5AttnProcessor,
]
def __init__(self, num_channels, head_dim, processor=None):
super().__init__()
assert num_channels % head_dim == 0
self.num_heads = num_channels // head_dim
self.to_query = nn.Linear(num_channels, num_channels, bias=True)
self.to_key = nn.Linear(num_channels, num_channels, bias=True)
self.to_value = nn.Linear(num_channels, num_channels, bias=True)
self.query_norm = nn.RMSNorm(head_dim)
self.key_norm = nn.RMSNorm(head_dim)
self.out_layer = nn.Linear(num_channels, num_channels, bias=True)
if processor is None:
processor = self._default_processor_cls()
self.set_processor(processor)
def forward(
self,
hidden_states: torch.Tensor,
encoder_hidden_states: Optional[torch.Tensor] = None,
sparse_params: Optional[torch.Tensor] = None,
rotary_emb: Optional[Tuple[torch.Tensor, torch.Tensor]] = None,
**kwargs,
) -> torch.Tensor:
attn_parameters = set(inspect.signature(self.processor.__call__).parameters.keys())
quiet_attn_parameters = {}
unused_kwargs = [k for k, _ in kwargs.items() if k not in attn_parameters and k not in quiet_attn_parameters]
if len(unused_kwargs) > 0:
logger.warning(
f"attention_processor_kwargs {unused_kwargs} are not expected by {self.processor.__class__.__name__} and will be ignored."
)
kwargs = {k: w for k, w in kwargs.items() if k in attn_parameters}
return self.processor(
self,
hidden_states,
encoder_hidden_states=encoder_hidden_states,
sparse_params=sparse_params,
rotary_emb=rotary_emb,
**kwargs,
)
class Kandinsky5FeedForward(nn.Module):
def __init__(self, dim, ff_dim):
super().__init__()
self.in_layer = nn.Linear(dim, ff_dim, bias=False)
self.activation = nn.GELU()
self.out_layer = nn.Linear(ff_dim, dim, bias=False)
def forward(self, x):
return self.out_layer(self.activation(self.in_layer(x)))
class Kandinsky5OutLayer(nn.Module):
def __init__(self, model_dim, time_dim, visual_dim, patch_size):
super().__init__()
self.patch_size = patch_size
self.modulation = Kandinsky5Modulation(time_dim, model_dim, 2)
self.norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.out_layer = nn.Linear(model_dim, math.prod(patch_size) * visual_dim, bias=True)
def forward(self, visual_embed, text_embed, time_embed):
shift, scale = torch.chunk(self.modulation(time_embed).unsqueeze(dim=1), 2, dim=-1)
visual_embed = (
self.norm(visual_embed.float()) * (scale.float()[:, None, None] + 1.0) + shift.float()[:, None, None]
).type_as(visual_embed)
x = self.out_layer(visual_embed)
batch_size, duration, height, width, _ = x.shape
x = (
x.view(
batch_size,
duration,
height,
width,
-1,
self.patch_size[0],
self.patch_size[1],
self.patch_size[2],
)
.permute(0, 1, 5, 2, 6, 3, 7, 4)
.flatten(1, 2)
.flatten(2, 3)
.flatten(3, 4)
)
return x
class Kandinsky5TransformerEncoderBlock(nn.Module):
def __init__(self, model_dim, time_dim, ff_dim, head_dim):
super().__init__()
self.text_modulation = Kandinsky5Modulation(time_dim, model_dim, 6)
self.self_attention_norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.self_attention = Kandinsky5Attention(model_dim, head_dim, processor=Kandinsky5AttnProcessor())
self.feed_forward_norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.feed_forward = Kandinsky5FeedForward(model_dim, ff_dim)
def forward(self, x, time_embed, rope):
self_attn_params, ff_params = torch.chunk(self.text_modulation(time_embed).unsqueeze(dim=1), 2, dim=-1)
shift, scale, gate = torch.chunk(self_attn_params, 3, dim=-1)
out = (self.self_attention_norm(x.float()) * (scale.float() + 1.0) + shift.float()).type_as(x)
out = self.self_attention(out, rotary_emb=rope)
x = (x.float() + gate.float() * out.float()).type_as(x)
shift, scale, gate = torch.chunk(ff_params, 3, dim=-1)
out = (self.feed_forward_norm(x.float()) * (scale.float() + 1.0) + shift.float()).type_as(x)
out = self.feed_forward(out)
x = (x.float() + gate.float() * out.float()).type_as(x)
return x
class Kandinsky5TransformerDecoderBlock(nn.Module):
def __init__(self, model_dim, time_dim, ff_dim, head_dim):
super().__init__()
self.visual_modulation = Kandinsky5Modulation(time_dim, model_dim, 9)
self.self_attention_norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.self_attention = Kandinsky5Attention(model_dim, head_dim, processor=Kandinsky5AttnProcessor())
self.cross_attention_norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.cross_attention = Kandinsky5Attention(model_dim, head_dim, processor=Kandinsky5AttnProcessor())
self.feed_forward_norm = nn.LayerNorm(model_dim, elementwise_affine=False)
self.feed_forward = Kandinsky5FeedForward(model_dim, ff_dim)
def forward(self, visual_embed, text_embed, time_embed, rope, sparse_params):
self_attn_params, cross_attn_params, ff_params = torch.chunk(
self.visual_modulation(time_embed).unsqueeze(dim=1), 3, dim=-1
)
shift, scale, gate = torch.chunk(self_attn_params, 3, dim=-1)
visual_out = (self.self_attention_norm(visual_embed.float()) * (scale.float() + 1.0) + shift.float()).type_as(
visual_embed
)
visual_out = self.self_attention(visual_out, rotary_emb=rope, sparse_params=sparse_params)
visual_embed = (visual_embed.float() + gate.float() * visual_out.float()).type_as(visual_embed)
shift, scale, gate = torch.chunk(cross_attn_params, 3, dim=-1)
visual_out = (self.cross_attention_norm(visual_embed.float()) * (scale.float() + 1.0) + shift.float()).type_as(
visual_embed
)
visual_out = self.cross_attention(visual_out, encoder_hidden_states=text_embed)
visual_embed = (visual_embed.float() + gate.float() * visual_out.float()).type_as(visual_embed)
shift, scale, gate = torch.chunk(ff_params, 3, dim=-1)
visual_out = (self.feed_forward_norm(visual_embed.float()) * (scale.float() + 1.0) + shift.float()).type_as(
visual_embed
)
visual_out = self.feed_forward(visual_out)
visual_embed = (visual_embed.float() + gate.float() * visual_out.float()).type_as(visual_embed)
return visual_embed
class Kandinsky5Transformer3DModel(
ModelMixin,
ConfigMixin,
PeftAdapterMixin,
FromOriginalModelMixin,
CacheMixin,
AttentionMixin,
):
"""
A 3D Diffusion Transformer model for video-like data.
"""
_repeated_blocks = [
"Kandinsky5TransformerEncoderBlock",
"Kandinsky5TransformerDecoderBlock",
]
_supports_gradient_checkpointing = True
@register_to_config
def __init__(
self,
in_visual_dim=4,
in_text_dim=3584,
in_text_dim2=768,
time_dim=512,
out_visual_dim=4,
patch_size=(1, 2, 2),
model_dim=2048,
ff_dim=5120,
num_text_blocks=2,
num_visual_blocks=32,
axes_dims=(16, 24, 24),
visual_cond=False,
attention_type: str = "regular",
attention_causal: bool = None,
attention_local: bool = None,
attention_glob: bool = None,
attention_window: int = None,
attention_P: float = None,
attention_wT: int = None,
attention_wW: int = None,
attention_wH: int = None,
attention_add_sta: bool = None,
attention_method: str = None,
):
super().__init__()
head_dim = sum(axes_dims)
self.in_visual_dim = in_visual_dim
self.model_dim = model_dim
self.patch_size = patch_size
self.visual_cond = visual_cond
self.attention_type = attention_type
visual_embed_dim = 2 * in_visual_dim + 1 if visual_cond else in_visual_dim
# Initialize embeddings
self.time_embeddings = Kandinsky5TimeEmbeddings(model_dim, time_dim)
self.text_embeddings = Kandinsky5TextEmbeddings(in_text_dim, model_dim)
self.pooled_text_embeddings = Kandinsky5TextEmbeddings(in_text_dim2, time_dim)
self.visual_embeddings = Kandinsky5VisualEmbeddings(visual_embed_dim, model_dim, patch_size)
# Initialize positional embeddings
self.text_rope_embeddings = Kandinsky5RoPE1D(head_dim)
self.visual_rope_embeddings = Kandinsky5RoPE3D(axes_dims)
# Initialize transformer blocks
self.text_transformer_blocks = nn.ModuleList(
[Kandinsky5TransformerEncoderBlock(model_dim, time_dim, ff_dim, head_dim) for _ in range(num_text_blocks)]
)
self.visual_transformer_blocks = nn.ModuleList(
[
Kandinsky5TransformerDecoderBlock(model_dim, time_dim, ff_dim, head_dim)
for _ in range(num_visual_blocks)
]
)
# Initialize output layer
self.out_layer = Kandinsky5OutLayer(model_dim, time_dim, out_visual_dim, patch_size)
self.gradient_checkpointing = False
def forward(
self,
hidden_states: torch.Tensor, # x
encoder_hidden_states: torch.Tensor, # text_embed
timestep: torch.Tensor, # time
pooled_projections: torch.Tensor, # pooled_text_embed
visual_rope_pos: Tuple[int, int, int],
text_rope_pos: torch.LongTensor,
scale_factor: Tuple[float, float, float] = (1.0, 1.0, 1.0),
sparse_params: Optional[Dict[str, Any]] = None,
return_dict: bool = True,
) -> Union[Transformer2DModelOutput, torch.FloatTensor]:
"""
Forward pass of the Kandinsky5 3D Transformer.
Args:
hidden_states (`torch.FloatTensor`): Input visual states
encoder_hidden_states (`torch.FloatTensor`): Text embeddings
timestep (`torch.Tensor` or `float` or `int`): Current timestep
pooled_projections (`torch.FloatTensor`): Pooled text embeddings
visual_rope_pos (`Tuple[int, int, int]`): Position for visual RoPE
text_rope_pos (`torch.LongTensor`): Position for text RoPE
scale_factor (`Tuple[float, float, float]`, optional): Scale factor for RoPE
sparse_params (`Dict[str, Any]`, optional): Parameters for sparse attention
return_dict (`bool`, optional): Whether to return a dictionary
Returns:
[`~models.transformer_2d.Transformer2DModelOutput`] or `torch.FloatTensor`: The output of the transformer
"""
x = hidden_states
text_embed = encoder_hidden_states
time = timestep
pooled_text_embed = pooled_projections
text_embed = self.text_embeddings(text_embed)
time_embed = self.time_embeddings(time)
time_embed = time_embed + self.pooled_text_embeddings(pooled_text_embed)
visual_embed = self.visual_embeddings(x)
text_rope = self.text_rope_embeddings(text_rope_pos)
text_rope = text_rope.unsqueeze(dim=0)
for text_transformer_block in self.text_transformer_blocks:
if torch.is_grad_enabled() and self.gradient_checkpointing:
text_embed = self._gradient_checkpointing_func(
text_transformer_block, text_embed, time_embed, text_rope
)
else:
text_embed = text_transformer_block(text_embed, time_embed, text_rope)
visual_shape = visual_embed.shape[:-1]
visual_rope = self.visual_rope_embeddings(visual_shape, visual_rope_pos, scale_factor)
to_fractal = sparse_params["to_fractal"] if sparse_params is not None else False
visual_embed, visual_rope = fractal_flatten(visual_embed, visual_rope, visual_shape, block_mask=to_fractal)
for visual_transformer_block in self.visual_transformer_blocks:
if torch.is_grad_enabled() and self.gradient_checkpointing:
visual_embed = self._gradient_checkpointing_func(
visual_transformer_block,
visual_embed,
text_embed,
time_embed,
visual_rope,
sparse_params,
)
else:
visual_embed = visual_transformer_block(
visual_embed, text_embed, time_embed, visual_rope, sparse_params
)
visual_embed = fractal_unflatten(visual_embed, visual_shape, block_mask=to_fractal)
x = self.out_layer(visual_embed, text_embed, time_embed)
if not return_dict:
return x
return Transformer2DModelOutput(sample=x)
@@ -180,7 +180,6 @@ class QwenEmbedRope(nn.Module):
],
dim=1,
)
self.rope_cache = {}
# DO NOT USING REGISTER BUFFER HERE, IT WILL CAUSE COMPLEX NUMBERS LOSE ITS IMAGINARY PART
self.scale_rope = scale_rope
@@ -195,10 +194,20 @@ class QwenEmbedRope(nn.Module):
freqs = torch.polar(torch.ones_like(freqs), freqs)
return freqs
def forward(self, video_fhw, txt_seq_lens, device):
def forward(
self,
video_fhw: Union[Tuple[int, int, int], List[Tuple[int, int, int]]],
txt_seq_lens: List[int],
device: torch.device,
) -> Tuple[torch.Tensor, torch.Tensor]:
"""
Args: video_fhw: [frame, height, width] a list of 3 integers representing the shape of the video Args:
txt_length: [bs] a list of 1 integers representing the length of the text
Args:
video_fhw (`Tuple[int, int, int]` or `List[Tuple[int, int, int]]`):
A list of 3 integers [frame, height, width] representing the shape of the video.
txt_seq_lens (`List[int]`):
A list of integers of length batch_size representing the length of each text prompt.
device: (`torch.device`):
The device on which to perform the RoPE computation.
"""
if self.pos_freqs.device != device:
self.pos_freqs = self.pos_freqs.to(device)
@@ -213,14 +222,8 @@ class QwenEmbedRope(nn.Module):
max_vid_index = 0
for idx, fhw in enumerate(video_fhw):
frame, height, width = fhw
rope_key = f"{idx}_{height}_{width}"
if not torch.compiler.is_compiling():
if rope_key not in self.rope_cache:
self.rope_cache[rope_key] = self._compute_video_freqs(frame, height, width, idx)
video_freq = self.rope_cache[rope_key]
else:
video_freq = self._compute_video_freqs(frame, height, width, idx)
# RoPE frequencies are cached via a lru_cache decorator on _compute_video_freqs
video_freq = self._compute_video_freqs(frame, height, width, idx)
video_freq = video_freq.to(device)
vid_freqs.append(video_freq)
@@ -235,8 +238,8 @@ class QwenEmbedRope(nn.Module):
return vid_freqs, txt_freqs
@functools.lru_cache(maxsize=None)
def _compute_video_freqs(self, frame, height, width, idx=0):
@functools.lru_cache(maxsize=128)
def _compute_video_freqs(self, frame: int, height: int, width: int, idx: int = 0) -> torch.Tensor:
seq_lens = frame * height * width
freqs_pos = self.pos_freqs.split([x // 2 for x in self.axes_dim], dim=1)
freqs_neg = self.neg_freqs.split([x // 2 for x in self.axes_dim], dim=1)
@@ -115,7 +115,7 @@ class StableDiffusionXLInpaintLoopBeforeDenoiser(ModularPipelineBlocks):
def check_inputs(components, block_state):
num_channels_unet = components.num_channels_unet
if num_channels_unet == 9:
# default case for runwayml/stable-diffusion-inpainting
# default case for stable-diffusion-v1-5/stable-diffusion-inpainting
if block_state.mask is None or block_state.masked_image_latents is None:
raise ValueError("mask and masked_image_latents must be provided for inpainting-specific Unet")
num_channels_latents = block_state.latents.shape[1]
+1 -1
View File
@@ -159,7 +159,7 @@ init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
"stable-diffusion-v1-5/stable-diffusion-inpainting",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
+2
View File
@@ -382,6 +382,7 @@ else:
"WuerstchenPriorPipeline",
]
_import_structure["wan"] = ["WanPipeline", "WanImageToVideoPipeline", "WanVideoToVideoPipeline", "WanVACEPipeline"]
_import_structure["kandinsky5"] = ["Kandinsky5T2VPipeline"]
_import_structure["skyreels_v2"] = [
"SkyReelsV2DiffusionForcingPipeline",
"SkyReelsV2DiffusionForcingImageToVideoPipeline",
@@ -671,6 +672,7 @@ if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
Kandinsky3Img2ImgPipeline,
Kandinsky3Pipeline,
)
from .kandinsky5 import Kandinsky5T2VPipeline
from .latent_consistency_models import (
LatentConsistencyModelImg2ImgPipeline,
LatentConsistencyModelPipeline,
@@ -133,8 +133,8 @@ class StableDiffusionControlNetXSPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -185,8 +185,8 @@ class AltDiffusionPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -266,8 +266,8 @@ class AltDiffusionPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -213,8 +213,8 @@ class AltDiffusionImg2ImgPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -294,8 +294,8 @@ class AltDiffusionImg2ImgPipeline(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -162,8 +162,8 @@ class CycleDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Sta
instance of [`DDIMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -226,8 +226,8 @@ class CycleDiffusionPipeline(DiffusionPipeline, TextualInversionLoaderMixin, Sta
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -62,7 +62,8 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -111,7 +111,8 @@ class StableDiffusionInpaintPipelineLegacy(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
@@ -196,8 +197,8 @@ class StableDiffusionInpaintPipelineLegacy(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
@@ -64,8 +64,8 @@ class StableDiffusionModelEditingPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
with_to_k ([`bool`]):
@@ -46,10 +46,12 @@ EXAMPLE_DOC_STRING = """
>>> from diffusers import DDPMParallelScheduler
>>> from diffusers import StableDiffusionParadigmsPipeline
>>> scheduler = DDPMParallelScheduler.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
>>> scheduler = DDPMParallelScheduler.from_pretrained(
... "stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="scheduler"
... )
>>> pipe = StableDiffusionParadigmsPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5", scheduler=scheduler, torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", scheduler=scheduler, torch_dtype=torch.float16
... )
>>> pipe = pipe.to("cuda")
@@ -95,8 +97,8 @@ class StableDiffusionParadigmsPipeline(
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -303,7 +303,8 @@ class StableDiffusionPix2PixZeroPipeline(DiffusionPipeline, StableDiffusionMixin
[`DDIMScheduler`], [`LMSDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`], or [`DDPMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
Please, refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
details.
feature_extractor ([`CLIPImageProcessor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
requires_safety_checker (bool):
@@ -38,8 +38,8 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
"""
@@ -0,0 +1,48 @@
from typing import TYPE_CHECKING
from ...utils import (
DIFFUSERS_SLOW_IMPORT,
OptionalDependencyNotAvailable,
_LazyModule,
get_objects_from_module,
is_torch_available,
is_transformers_available,
)
_dummy_objects = {}
_import_structure = {}
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils import dummy_torch_and_transformers_objects # noqa F403
_dummy_objects.update(get_objects_from_module(dummy_torch_and_transformers_objects))
else:
_import_structure["pipeline_kandinsky"] = ["Kandinsky5T2VPipeline"]
if TYPE_CHECKING or DIFFUSERS_SLOW_IMPORT:
try:
if not (is_transformers_available() and is_torch_available()):
raise OptionalDependencyNotAvailable()
except OptionalDependencyNotAvailable:
from ...utils.dummy_torch_and_transformers_objects import *
else:
from .pipeline_kandinsky import Kandinsky5T2VPipeline
else:
import sys
sys.modules[__name__] = _LazyModule(
__name__,
globals()["__file__"],
_import_structure,
module_spec=__spec__,
)
for name, value in _dummy_objects.items():
setattr(sys.modules[__name__], name, value)
@@ -0,0 +1,893 @@
# Copyright 2025 The Kandinsky Team and The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import html
from typing import Callable, Dict, List, Optional, Union
import regex as re
import torch
from torch.nn import functional as F
from transformers import CLIPTextModel, CLIPTokenizer, Qwen2_5_VLForConditionalGeneration, Qwen2VLProcessor
from ...callbacks import MultiPipelineCallbacks, PipelineCallback
from ...loaders import KandinskyLoraLoaderMixin
from ...models import AutoencoderKLHunyuanVideo
from ...models.transformers import Kandinsky5Transformer3DModel
from ...schedulers import FlowMatchEulerDiscreteScheduler
from ...utils import is_ftfy_available, is_torch_xla_available, logging, replace_example_docstring
from ...utils.torch_utils import randn_tensor
from ...video_processor import VideoProcessor
from ..pipeline_utils import DiffusionPipeline
from .pipeline_output import KandinskyPipelineOutput
if is_torch_xla_available():
import torch_xla.core.xla_model as xm
XLA_AVAILABLE = True
else:
XLA_AVAILABLE = False
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
if is_ftfy_available():
import ftfy
logger = logging.get_logger(__name__)
EXAMPLE_DOC_STRING = """
Examples:
```python
>>> import torch
>>> from diffusers import Kandinsky5T2VPipeline
>>> from diffusers.utils import export_to_video
>>> # Available models:
>>> # ai-forever/Kandinsky-5.0-T2V-Lite-sft-5s-Diffusers
>>> # ai-forever/Kandinsky-5.0-T2V-Lite-nocfg-5s-Diffusers
>>> # ai-forever/Kandinsky-5.0-T2V-Lite-distilled16steps-5s-Diffusers
>>> # ai-forever/Kandinsky-5.0-T2V-Lite-pretrain-5s-Diffusers
>>> model_id = "ai-forever/Kandinsky-5.0-T2V-Lite-sft-5s-Diffusers"
>>> pipe = Kandinsky5T2VPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
>>> pipe = pipe.to("cuda")
>>> prompt = "A cat and a dog baking a cake together in a kitchen."
>>> negative_prompt = "Static, 2D cartoon, cartoon, 2d animation, paintings, images, worst quality, low quality, ugly, deformed, walking backwards"
>>> output = pipe(
... prompt=prompt,
... negative_prompt=negative_prompt,
... height=512,
... width=768,
... num_frames=121,
... num_inference_steps=50,
... guidance_scale=5.0,
... ).frames[0]
>>> export_to_video(output, "output.mp4", fps=24, quality=9)
```
"""
def basic_clean(text):
"""Clean text using ftfy if available and unescape HTML entities."""
if is_ftfy_available():
text = ftfy.fix_text(text)
text = html.unescape(html.unescape(text))
return text.strip()
def whitespace_clean(text):
"""Normalize whitespace in text by replacing multiple spaces with single space."""
text = re.sub(r"\s+", " ", text)
text = text.strip()
return text
def prompt_clean(text):
"""Apply both basic cleaning and whitespace normalization to prompts."""
text = whitespace_clean(basic_clean(text))
return text
class Kandinsky5T2VPipeline(DiffusionPipeline, KandinskyLoraLoaderMixin):
r"""
Pipeline for text-to-video generation using Kandinsky 5.0.
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
Args:
transformer ([`Kandinsky5Transformer3DModel`]):
Conditional Transformer to denoise the encoded video latents.
vae ([`AutoencoderKLHunyuanVideo`]):
Variational Auto-Encoder (VAE) Model to encode and decode videos to and from latent representations.
text_encoder ([`Qwen2_5_VLForConditionalGeneration`]):
Frozen text-encoder (Qwen2.5-VL).
tokenizer ([`AutoProcessor`]):
Tokenizer for Qwen2.5-VL.
text_encoder_2 ([`CLIPTextModel`]):
Frozen CLIP text encoder.
tokenizer_2 ([`CLIPTokenizer`]):
Tokenizer for CLIP.
scheduler ([`FlowMatchEulerDiscreteScheduler`]):
A scheduler to be used in combination with `transformer` to denoise the encoded video latents.
"""
model_cpu_offload_seq = "text_encoder->text_encoder_2->transformer->vae"
_callback_tensor_inputs = [
"latents",
"prompt_embeds_qwen",
"prompt_embeds_clip",
"negative_prompt_embeds_qwen",
"negative_prompt_embeds_clip",
]
def __init__(
self,
transformer: Kandinsky5Transformer3DModel,
vae: AutoencoderKLHunyuanVideo,
text_encoder: Qwen2_5_VLForConditionalGeneration,
tokenizer: Qwen2VLProcessor,
text_encoder_2: CLIPTextModel,
tokenizer_2: CLIPTokenizer,
scheduler: FlowMatchEulerDiscreteScheduler,
):
super().__init__()
self.register_modules(
transformer=transformer,
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
text_encoder_2=text_encoder_2,
tokenizer_2=tokenizer_2,
scheduler=scheduler,
)
self.prompt_template = "\n".join(
[
"<|im_start|>system\nYou are a promt engineer. Describe the video in detail.",
"Describe how the camera moves or shakes, describe the zoom and view angle, whether it follows the objects.",
"Describe the location of the video, main characters or objects and their action.",
"Describe the dynamism of the video and presented actions.",
"Name the visual style of the video: whether it is a professional footage, user generated content, some kind of animation, video game or scren content.",
"Describe the visual effects, postprocessing and transitions if they are presented in the video.",
"Pay attention to the order of key actions shown in the scene.<|im_end|>",
"<|im_start|>user\n{}<|im_end|>",
]
)
self.prompt_template_encode_start_idx = 129
self.vae_scale_factor_temporal = vae.config.temporal_compression_ratio
self.vae_scale_factor_spatial = vae.config.spatial_compression_ratio
self.video_processor = VideoProcessor(vae_scale_factor=self.vae_scale_factor_spatial)
@staticmethod
def fast_sta_nabla(T: int, H: int, W: int, wT: int = 3, wH: int = 3, wW: int = 3, device="cuda") -> torch.Tensor:
"""
Create a sparse temporal attention (STA) mask for efficient video generation.
This method generates a mask that limits attention to nearby frames and spatial positions, reducing
computational complexity for video generation.
Args:
T (int): Number of temporal frames
H (int): Height in latent space
W (int): Width in latent space
wT (int): Temporal attention window size
wH (int): Height attention window size
wW (int): Width attention window size
device (str): Device to create tensor on
Returns:
torch.Tensor: Sparse attention mask of shape (T*H*W, T*H*W)
"""
l = torch.Tensor([T, H, W]).amax()
r = torch.arange(0, l, 1, dtype=torch.int16, device=device)
mat = (r.unsqueeze(1) - r.unsqueeze(0)).abs()
sta_t, sta_h, sta_w = (
mat[:T, :T].flatten(),
mat[:H, :H].flatten(),
mat[:W, :W].flatten(),
)
sta_t = sta_t <= wT // 2
sta_h = sta_h <= wH // 2
sta_w = sta_w <= wW // 2
sta_hw = (sta_h.unsqueeze(1) * sta_w.unsqueeze(0)).reshape(H, H, W, W).transpose(1, 2).flatten()
sta = (sta_t.unsqueeze(1) * sta_hw.unsqueeze(0)).reshape(T, T, H * W, H * W).transpose(1, 2)
return sta.reshape(T * H * W, T * H * W)
def get_sparse_params(self, sample, device):
"""
Generate sparse attention parameters for the transformer based on sample dimensions.
This method computes the sparse attention configuration needed for efficient video processing in the
transformer model.
Args:
sample (torch.Tensor): Input sample tensor
device (torch.device): Device to place tensors on
Returns:
Dict: Dictionary containing sparse attention parameters
"""
assert self.transformer.config.patch_size[0] == 1
B, T, H, W, _ = sample.shape
T, H, W = (
T // self.transformer.config.patch_size[0],
H // self.transformer.config.patch_size[1],
W // self.transformer.config.patch_size[2],
)
if self.transformer.config.attention_type == "nabla":
sta_mask = self.fast_sta_nabla(
T,
H // 8,
W // 8,
self.transformer.config.attention_wT,
self.transformer.config.attention_wH,
self.transformer.config.attention_wW,
device=device,
)
sparse_params = {
"sta_mask": sta_mask.unsqueeze_(0).unsqueeze_(0),
"attention_type": self.transformer.config.attention_type,
"to_fractal": True,
"P": self.transformer.config.attention_P,
"wT": self.transformer.config.attention_wT,
"wW": self.transformer.config.attention_wW,
"wH": self.transformer.config.attention_wH,
"add_sta": self.transformer.config.attention_add_sta,
"visual_shape": (T, H, W),
"method": self.transformer.config.attention_method,
}
else:
sparse_params = None
return sparse_params
def _encode_prompt_qwen(
self,
prompt: Union[str, List[str]],
device: Optional[torch.device] = None,
max_sequence_length: int = 256,
dtype: Optional[torch.dtype] = None,
):
"""
Encode prompt using Qwen2.5-VL text encoder.
This method processes the input prompt through the Qwen2.5-VL model to generate text embeddings suitable for
video generation.
Args:
prompt (Union[str, List[str]]): Input prompt or list of prompts
device (torch.device): Device to run encoding on
num_videos_per_prompt (int): Number of videos to generate per prompt
max_sequence_length (int): Maximum sequence length for tokenization
dtype (torch.dtype): Data type for embeddings
Returns:
Tuple[torch.Tensor, torch.Tensor]: Text embeddings and cumulative sequence lengths
"""
device = device or self._execution_device
dtype = dtype or self.text_encoder.dtype
full_texts = [self.prompt_template.format(p) for p in prompt]
inputs = self.tokenizer(
text=full_texts,
images=None,
videos=None,
max_length=max_sequence_length + self.prompt_template_encode_start_idx,
truncation=True,
return_tensors="pt",
padding=True,
).to(device)
embeds = self.text_encoder(
input_ids=inputs["input_ids"],
return_dict=True,
output_hidden_states=True,
)["hidden_states"][-1][:, self.prompt_template_encode_start_idx :]
attention_mask = inputs["attention_mask"][:, self.prompt_template_encode_start_idx :]
cu_seqlens = torch.cumsum(attention_mask.sum(1), dim=0)
cu_seqlens = F.pad(cu_seqlens, (1, 0), value=0).to(dtype=torch.int32)
return embeds.to(dtype), cu_seqlens
def _encode_prompt_clip(
self,
prompt: Union[str, List[str]],
device: Optional[torch.device] = None,
dtype: Optional[torch.dtype] = None,
):
"""
Encode prompt using CLIP text encoder.
This method processes the input prompt through the CLIP model to generate pooled embeddings that capture
semantic information.
Args:
prompt (Union[str, List[str]]): Input prompt or list of prompts
device (torch.device): Device to run encoding on
num_videos_per_prompt (int): Number of videos to generate per prompt
dtype (torch.dtype): Data type for embeddings
Returns:
torch.Tensor: Pooled text embeddings from CLIP
"""
device = device or self._execution_device
dtype = dtype or self.text_encoder_2.dtype
inputs = self.tokenizer_2(
prompt,
max_length=77,
truncation=True,
add_special_tokens=True,
padding="max_length",
return_tensors="pt",
).to(device)
pooled_embed = self.text_encoder_2(**inputs)["pooler_output"]
return pooled_embed.to(dtype)
def encode_prompt(
self,
prompt: Union[str, List[str]],
num_videos_per_prompt: int = 1,
max_sequence_length: int = 512,
device: Optional[torch.device] = None,
dtype: Optional[torch.dtype] = None,
):
r"""
Encodes a single prompt (positive or negative) into text encoder hidden states.
This method combines embeddings from both Qwen2.5-VL and CLIP text encoders to create comprehensive text
representations for video generation.
Args:
prompt (`str` or `List[str]`):
Prompt to be encoded.
num_videos_per_prompt (`int`, *optional*, defaults to 1):
Number of videos to generate per prompt.
max_sequence_length (`int`, *optional*, defaults to 512):
Maximum sequence length for text encoding.
device (`torch.device`, *optional*):
Torch device.
dtype (`torch.dtype`, *optional*):
Torch dtype.
Returns:
Tuple[torch.Tensor, torch.Tensor, torch.Tensor]:
- Qwen text embeddings of shape (batch_size * num_videos_per_prompt, sequence_length, embedding_dim)
- CLIP pooled embeddings of shape (batch_size * num_videos_per_prompt, clip_embedding_dim)
- Cumulative sequence lengths (`cu_seqlens`) for Qwen embeddings of shape (batch_size *
num_videos_per_prompt + 1,)
"""
device = device or self._execution_device
dtype = dtype or self.text_encoder.dtype
batch_size = len(prompt)
prompt = [prompt_clean(p) for p in prompt]
# Encode with Qwen2.5-VL
prompt_embeds_qwen, prompt_cu_seqlens = self._encode_prompt_qwen(
prompt=prompt,
device=device,
max_sequence_length=max_sequence_length,
dtype=dtype,
)
# prompt_embeds_qwen shape: [batch_size, seq_len, embed_dim]
# Encode with CLIP
prompt_embeds_clip = self._encode_prompt_clip(
prompt=prompt,
device=device,
dtype=dtype,
)
# prompt_embeds_clip shape: [batch_size, clip_embed_dim]
# Repeat embeddings for num_videos_per_prompt
# Qwen embeddings: repeat sequence for each video, then reshape
prompt_embeds_qwen = prompt_embeds_qwen.repeat(
1, num_videos_per_prompt, 1
) # [batch_size, seq_len * num_videos_per_prompt, embed_dim]
# Reshape to [batch_size * num_videos_per_prompt, seq_len, embed_dim]
prompt_embeds_qwen = prompt_embeds_qwen.view(
batch_size * num_videos_per_prompt, -1, prompt_embeds_qwen.shape[-1]
)
# CLIP embeddings: repeat for each video
prompt_embeds_clip = prompt_embeds_clip.repeat(
1, num_videos_per_prompt, 1
) # [batch_size, num_videos_per_prompt, clip_embed_dim]
# Reshape to [batch_size * num_videos_per_prompt, clip_embed_dim]
prompt_embeds_clip = prompt_embeds_clip.view(batch_size * num_videos_per_prompt, -1)
# Repeat cumulative sequence lengths for num_videos_per_prompt
# Original cu_seqlens: [0, len1, len1+len2, ...]
# Need to repeat the differences and reconstruct for repeated prompts
# Original differences (lengths) for each prompt in the batch
original_lengths = prompt_cu_seqlens.diff() # [len1, len2, ...]
# Repeat the lengths for num_videos_per_prompt
repeated_lengths = original_lengths.repeat_interleave(
num_videos_per_prompt
) # [len1, len1, ..., len2, len2, ...]
# Reconstruct the cumulative lengths
repeated_cu_seqlens = torch.cat(
[torch.tensor([0], device=device, dtype=torch.int32), repeated_lengths.cumsum(0)]
)
return prompt_embeds_qwen, prompt_embeds_clip, repeated_cu_seqlens
def check_inputs(
self,
prompt,
negative_prompt,
height,
width,
prompt_embeds_qwen=None,
prompt_embeds_clip=None,
negative_prompt_embeds_qwen=None,
negative_prompt_embeds_clip=None,
prompt_cu_seqlens=None,
negative_prompt_cu_seqlens=None,
callback_on_step_end_tensor_inputs=None,
):
"""
Validate input parameters for the pipeline.
Args:
prompt: Input prompt
negative_prompt: Negative prompt for guidance
height: Video height
width: Video width
prompt_embeds_qwen: Pre-computed Qwen prompt embeddings
prompt_embeds_clip: Pre-computed CLIP prompt embeddings
negative_prompt_embeds_qwen: Pre-computed Qwen negative prompt embeddings
negative_prompt_embeds_clip: Pre-computed CLIP negative prompt embeddings
prompt_cu_seqlens: Pre-computed cumulative sequence lengths for Qwen positive prompt
negative_prompt_cu_seqlens: Pre-computed cumulative sequence lengths for Qwen negative prompt
callback_on_step_end_tensor_inputs: Callback tensor inputs
Raises:
ValueError: If inputs are invalid
"""
if height % 16 != 0 or width % 16 != 0:
raise ValueError(f"`height` and `width` have to be divisible by 16 but are {height} and {width}.")
if callback_on_step_end_tensor_inputs is not None and not all(
k in self._callback_tensor_inputs for k in callback_on_step_end_tensor_inputs
):
raise ValueError(
f"`callback_on_step_end_tensor_inputs` has to be in {self._callback_tensor_inputs}, but found {[k for k in callback_on_step_end_tensor_inputs if k not in self._callback_tensor_inputs]}"
)
# Check for consistency within positive prompt embeddings and sequence lengths
if prompt_embeds_qwen is not None or prompt_embeds_clip is not None or prompt_cu_seqlens is not None:
if prompt_embeds_qwen is None or prompt_embeds_clip is None or prompt_cu_seqlens is None:
raise ValueError(
"If any of `prompt_embeds_qwen`, `prompt_embeds_clip`, or `prompt_cu_seqlens` is provided, "
"all three must be provided."
)
# Check for consistency within negative prompt embeddings and sequence lengths
if (
negative_prompt_embeds_qwen is not None
or negative_prompt_embeds_clip is not None
or negative_prompt_cu_seqlens is not None
):
if (
negative_prompt_embeds_qwen is None
or negative_prompt_embeds_clip is None
or negative_prompt_cu_seqlens is None
):
raise ValueError(
"If any of `negative_prompt_embeds_qwen`, `negative_prompt_embeds_clip`, or `negative_prompt_cu_seqlens` is provided, "
"all three must be provided."
)
# Check if prompt or embeddings are provided (either prompt or all required embedding components for positive)
if prompt is None and prompt_embeds_qwen is None:
raise ValueError(
"Provide either `prompt` or `prompt_embeds_qwen` (and corresponding `prompt_embeds_clip` and `prompt_cu_seqlens`). Cannot leave all undefined."
)
# Validate types for prompt and negative_prompt if provided
if prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if negative_prompt is not None and (
not isinstance(negative_prompt, str) and not isinstance(negative_prompt, list)
):
raise ValueError(f"`negative_prompt` has to be of type `str` or `list` but is {type(negative_prompt)}")
def prepare_latents(
self,
batch_size: int,
num_channels_latents: int = 16,
height: int = 480,
width: int = 832,
num_frames: int = 81,
dtype: Optional[torch.dtype] = None,
device: Optional[torch.device] = None,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
latents: Optional[torch.Tensor] = None,
) -> torch.Tensor:
"""
Prepare initial latent variables for video generation.
This method creates random noise latents or uses provided latents as starting point for the denoising process.
Args:
batch_size (int): Number of videos to generate
num_channels_latents (int): Number of channels in latent space
height (int): Height of generated video
width (int): Width of generated video
num_frames (int): Number of frames in video
dtype (torch.dtype): Data type for latents
device (torch.device): Device to create latents on
generator (torch.Generator): Random number generator
latents (torch.Tensor): Pre-existing latents to use
Returns:
torch.Tensor: Prepared latent tensor
"""
if latents is not None:
return latents.to(device=device, dtype=dtype)
num_latent_frames = (num_frames - 1) // self.vae_scale_factor_temporal + 1
shape = (
batch_size,
num_latent_frames,
int(height) // self.vae_scale_factor_spatial,
int(width) // self.vae_scale_factor_spatial,
num_channels_latents,
)
if isinstance(generator, list) and len(generator) != batch_size:
raise ValueError(
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch"
f" size of {batch_size}. Make sure the batch size matches the length of the generators."
)
latents = randn_tensor(shape, generator=generator, device=device, dtype=dtype)
if self.transformer.visual_cond:
# For visual conditioning, concatenate with zeros and mask
visual_cond = torch.zeros_like(latents)
visual_cond_mask = torch.zeros(
[
batch_size,
num_latent_frames,
int(height) // self.vae_scale_factor_spatial,
int(width) // self.vae_scale_factor_spatial,
1,
],
dtype=latents.dtype,
device=latents.device,
)
latents = torch.cat([latents, visual_cond, visual_cond_mask], dim=-1)
return latents
@property
def guidance_scale(self):
"""Get the current guidance scale value."""
return self._guidance_scale
@property
def do_classifier_free_guidance(self):
"""Check if classifier-free guidance is enabled."""
return self._guidance_scale > 1.0
@property
def num_timesteps(self):
"""Get the number of denoising timesteps."""
return self._num_timesteps
@property
def interrupt(self):
"""Check if generation has been interrupted."""
return self._interrupt
@torch.no_grad()
@replace_example_docstring(EXAMPLE_DOC_STRING)
def __call__(
self,
prompt: Union[str, List[str]] = None,
negative_prompt: Optional[Union[str, List[str]]] = None,
height: int = 512,
width: int = 768,
num_frames: int = 121,
num_inference_steps: int = 50,
guidance_scale: float = 5.0,
num_videos_per_prompt: Optional[int] = 1,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
latents: Optional[torch.Tensor] = None,
prompt_embeds_qwen: Optional[torch.Tensor] = None,
prompt_embeds_clip: Optional[torch.Tensor] = None,
negative_prompt_embeds_qwen: Optional[torch.Tensor] = None,
negative_prompt_embeds_clip: Optional[torch.Tensor] = None,
prompt_cu_seqlens: Optional[torch.Tensor] = None,
negative_prompt_cu_seqlens: Optional[torch.Tensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback_on_step_end: Optional[
Union[Callable[[int, int, Dict], None], PipelineCallback, MultiPipelineCallbacks]
] = None,
callback_on_step_end_tensor_inputs: List[str] = ["latents"],
max_sequence_length: int = 512,
**kwargs,
):
r"""
The call function to the pipeline for generation.
Args:
prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to guide the video generation. If not defined, pass `prompt_embeds` instead.
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts to avoid during video generation. If not defined, pass `negative_prompt_embeds`
instead. Ignored when not using guidance (`guidance_scale` < `1`).
height (`int`, defaults to `512`):
The height in pixels of the generated video.
width (`int`, defaults to `768`):
The width in pixels of the generated video.
num_frames (`int`, defaults to `25`):
The number of frames in the generated video.
num_inference_steps (`int`, defaults to `50`):
The number of denoising steps.
guidance_scale (`float`, defaults to `5.0`):
Guidance scale as defined in classifier-free guidance.
num_videos_per_prompt (`int`, *optional*, defaults to 1):
The number of videos to generate per prompt.
generator (`torch.Generator` or `List[torch.Generator]`, *optional*):
A torch generator to make generation deterministic.
latents (`torch.Tensor`, *optional*):
Pre-generated noisy latents.
prompt_embeds (`torch.Tensor`, *optional*):
Pre-generated text embeddings.
negative_prompt_embeds (`torch.Tensor`, *optional*):
Pre-generated negative text embeddings.
output_type (`str`, *optional*, defaults to `"pil"`):
The output format of the generated video.
return_dict (`bool`, *optional*, defaults to `True`):
Whether or not to return a [`KandinskyPipelineOutput`].
callback_on_step_end (`Callable`, `PipelineCallback`, `MultiPipelineCallbacks`, *optional*):
A function that is called at the end of each denoising step.
callback_on_step_end_tensor_inputs (`List`, *optional*):
The list of tensor inputs for the `callback_on_step_end` function.
max_sequence_length (`int`, defaults to `512`):
The maximum sequence length for text encoding.
Examples:
Returns:
[`~KandinskyPipelineOutput`] or `tuple`:
If `return_dict` is `True`, [`KandinskyPipelineOutput`] is returned, otherwise a `tuple` is returned
where the first element is a list with the generated images.
"""
if isinstance(callback_on_step_end, (PipelineCallback, MultiPipelineCallbacks)):
callback_on_step_end_tensor_inputs = callback_on_step_end.tensor_inputs
# 1. Check inputs. Raise error if not correct
self.check_inputs(
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
prompt_embeds_qwen=prompt_embeds_qwen,
prompt_embeds_clip=prompt_embeds_clip,
negative_prompt_embeds_qwen=negative_prompt_embeds_qwen,
negative_prompt_embeds_clip=negative_prompt_embeds_clip,
prompt_cu_seqlens=prompt_cu_seqlens,
negative_prompt_cu_seqlens=negative_prompt_cu_seqlens,
callback_on_step_end_tensor_inputs=callback_on_step_end_tensor_inputs,
)
if num_frames % self.vae_scale_factor_temporal != 1:
logger.warning(
f"`num_frames - 1` has to be divisible by {self.vae_scale_factor_temporal}. Rounding to the nearest number."
)
num_frames = num_frames // self.vae_scale_factor_temporal * self.vae_scale_factor_temporal + 1
num_frames = max(num_frames, 1)
self._guidance_scale = guidance_scale
self._interrupt = False
device = self._execution_device
dtype = self.transformer.dtype
# 2. Define call parameters
if prompt is not None and isinstance(prompt, str):
batch_size = 1
prompt = [prompt]
elif prompt is not None and isinstance(prompt, list):
batch_size = len(prompt)
else:
batch_size = prompt_embeds_qwen.shape[0]
# 3. Encode input prompt
if prompt_embeds_qwen is None:
prompt_embeds_qwen, prompt_embeds_clip, prompt_cu_seqlens = self.encode_prompt(
prompt=prompt,
max_sequence_length=max_sequence_length,
device=device,
dtype=dtype,
)
if self.do_classifier_free_guidance:
if negative_prompt is None:
negative_prompt = "Static, 2D cartoon, cartoon, 2d animation, paintings, images, worst quality, low quality, ugly, deformed, walking backwards"
if isinstance(negative_prompt, str):
negative_prompt = [negative_prompt] * len(prompt) if prompt is not None else [negative_prompt]
elif len(negative_prompt) != len(prompt):
raise ValueError(
f"`negative_prompt` must have same length as `prompt`. Got {len(negative_prompt)} vs {len(prompt)}."
)
if negative_prompt_embeds_qwen is None:
negative_prompt_embeds_qwen, negative_prompt_embeds_clip, negative_cu_seqlens = self.encode_prompt(
prompt=negative_prompt,
max_sequence_length=max_sequence_length,
device=device,
dtype=dtype,
)
# 4. Prepare timesteps
self.scheduler.set_timesteps(num_inference_steps, device=device)
timesteps = self.scheduler.timesteps
# 5. Prepare latent variables
num_channels_latents = self.transformer.config.in_visual_dim
latents = self.prepare_latents(
batch_size * num_videos_per_prompt,
num_channels_latents,
height,
width,
num_frames,
dtype,
device,
generator,
latents,
)
# 6. Prepare rope positions for positional encoding
num_latent_frames = (num_frames - 1) // self.vae_scale_factor_temporal + 1
visual_rope_pos = [
torch.arange(num_latent_frames, device=device),
torch.arange(height // self.vae_scale_factor_spatial // 2, device=device),
torch.arange(width // self.vae_scale_factor_spatial // 2, device=device),
]
text_rope_pos = torch.arange(prompt_cu_seqlens.diff().max().item(), device=device)
negative_text_rope_pos = (
torch.arange(negative_cu_seqlens.diff().max().item(), device=device)
if negative_cu_seqlens is not None
else None
)
# 7. Sparse Params for efficient attention
sparse_params = self.get_sparse_params(latents, device)
# 8. Denoising loop
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
self._num_timesteps = len(timesteps)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
if self.interrupt:
continue
timestep = t.unsqueeze(0).repeat(batch_size * num_videos_per_prompt)
# Predict noise residual
pred_velocity = self.transformer(
hidden_states=latents.to(dtype),
encoder_hidden_states=prompt_embeds_qwen.to(dtype),
pooled_projections=prompt_embeds_clip.to(dtype),
timestep=timestep.to(dtype),
visual_rope_pos=visual_rope_pos,
text_rope_pos=text_rope_pos,
scale_factor=(1, 2, 2),
sparse_params=sparse_params,
return_dict=True,
).sample
if self.do_classifier_free_guidance and negative_prompt_embeds_qwen is not None:
uncond_pred_velocity = self.transformer(
hidden_states=latents.to(dtype),
encoder_hidden_states=negative_prompt_embeds_qwen.to(dtype),
pooled_projections=negative_prompt_embeds_clip.to(dtype),
timestep=timestep.to(dtype),
visual_rope_pos=visual_rope_pos,
text_rope_pos=negative_text_rope_pos,
scale_factor=(1, 2, 2),
sparse_params=sparse_params,
return_dict=True,
).sample
pred_velocity = uncond_pred_velocity + guidance_scale * (pred_velocity - uncond_pred_velocity)
# Compute previous sample using the scheduler
latents[:, :, :, :, :num_channels_latents] = self.scheduler.step(
pred_velocity, t, latents[:, :, :, :, :num_channels_latents], return_dict=False
)[0]
if callback_on_step_end is not None:
callback_kwargs = {}
for k in callback_on_step_end_tensor_inputs:
callback_kwargs[k] = locals()[k]
callback_outputs = callback_on_step_end(self, i, t, callback_kwargs)
latents = callback_outputs.pop("latents", latents)
prompt_embeds_qwen = callback_outputs.pop("prompt_embeds_qwen", prompt_embeds_qwen)
prompt_embeds_clip = callback_outputs.pop("prompt_embeds_clip", prompt_embeds_clip)
negative_prompt_embeds_qwen = callback_outputs.pop(
"negative_prompt_embeds_qwen", negative_prompt_embeds_qwen
)
negative_prompt_embeds_clip = callback_outputs.pop(
"negative_prompt_embeds_clip", negative_prompt_embeds_clip
)
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0):
progress_bar.update()
if XLA_AVAILABLE:
xm.mark_step()
# 8. Post-processing - extract main latents
latents = latents[:, :, :, :, :num_channels_latents]
# 9. Decode latents to video
if output_type != "latent":
latents = latents.to(self.vae.dtype)
# Reshape and normalize latents
video = latents.reshape(
batch_size,
num_videos_per_prompt,
(num_frames - 1) // self.vae_scale_factor_temporal + 1,
height // self.vae_scale_factor_spatial,
width // self.vae_scale_factor_spatial,
num_channels_latents,
)
video = video.permute(0, 1, 5, 2, 3, 4) # [batch, num_videos, channels, frames, height, width]
video = video.reshape(
batch_size * num_videos_per_prompt,
num_channels_latents,
(num_frames - 1) // self.vae_scale_factor_temporal + 1,
height // self.vae_scale_factor_spatial,
width // self.vae_scale_factor_spatial,
)
# Normalize and decode through VAE
video = video / self.vae.config.scaling_factor
video = self.vae.decode(video).sample
video = self.video_processor.postprocess_video(video, output_type=output_type)
else:
video = latents
# Offload all models
self.maybe_free_model_hooks()
if not return_dict:
return (video,)
return KandinskyPipelineOutput(frames=video)
@@ -0,0 +1,20 @@
from dataclasses import dataclass
import torch
from diffusers.utils import BaseOutput
@dataclass
class KandinskyPipelineOutput(BaseOutput):
r"""
Output class for Wan pipelines.
Args:
frames (`torch.Tensor`, `np.ndarray`, or List[List[PIL.Image.Image]]):
List of video outputs - It can be a nested list of length `batch_size,` with each sub-list containing
denoised PIL image sequences of length `num_frames.` It can also be a NumPy array or Torch tensor of shape
`(batch_size, num_frames, channels, height, width)`.
"""
frames: torch.Tensor
@@ -186,8 +186,8 @@ class LatentConsistencyModelImg2ImgPipeline(
supports [`LCMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
requires_safety_checker (`bool`, *optional*, defaults to `True`):
@@ -165,8 +165,8 @@ class LatentConsistencyModelPipeline(
supports [`LCMScheduler`].
safety_checker ([`StableDiffusionSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for more details
about a model's potential harms.
Please refer to the [model card](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) for
more details about a model's potential harms.
feature_extractor ([`~transformers.CLIPImageProcessor`]):
A `CLIPImageProcessor` to extract features from generated images; used as inputs to the `safety_checker`.
requires_safety_checker (`bool`, *optional*, defaults to `True`):
@@ -49,7 +49,7 @@ EXAMPLE_DOC_STRING = """
>>> from diffusers.utils import load_image
>>> pipe = LEditsPPPipelineStableDiffusion.from_pretrained(
... "runwayml/stable-diffusion-v1-5", variant="fp16", torch_dtype=torch.float16
... "stable-diffusion-v1-5/stable-diffusion-v1-5", variant="fp16", torch_dtype=torch.float16
... )
>>> pipe.enable_vae_tiling()
>>> pipe = pipe.to("cuda")
@@ -381,8 +381,8 @@ class LEditsPPPipelineStableDiffusion(
"The configuration file of the unet has set the default `sample_size` to smaller than"
" 64 which seems highly unlikely. If your checkpoint is a fine-tuned version of any of the"
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- stable-diffusion-v1-5/stable-diffusion-v1-5"
" \n- stable-diffusion-v1-5/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
" in the config might lead to incorrect results in future versions. If you have downloaded this"
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"

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